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Abhilasha Shourie Shilpa S. Chapadgaonkar Bioanalytical techniques form an integral part of applied biology and biomedical sciences. The book provides understanding of the concept and working principles of various bioanalytical techniques used in biomedical sciences, environmental studies, life sciences, pharmaceutical analysis, molecular biology, and biotechnological research, as well as the various instruments used in these processes. Divided into 12 chapters, the book provides a comprehensive account of microscopy, centrifugation, chromatography, electrophoresis, spectroscopy. It also focuses on two main topics: radioisotope and immunodiagnostic techniques. Techniques in molecular biology and recombinant DNA technology have also been described in detail.
• Explains analytical instrumentation in a concise manner • Provides state-of-the-art sophisticated techniques that would be beneficial to researchers in various fields for experimentation • Encourages reader to analytical thinking and practical application of the technique
Bioanalytical Techniques Abhilasha Shourie Shilpa S. Chapadgaonkar
Shourie • Chapadgaonkar
Key Features
Bioanalytical Techniques
Bioanalytical Techniques
The Energy and Resources Institute
Abhilasha Shourie Shilpa S. Chapadgaonkar
The Energy and Resources Institute
© The Energy and Resources Institute, 2015
ISBN 978-81-7993-529-3
All rights reserved. No part of this publication may be reproduced, stored in a retrieval system, or transmitted in any form or by any means, electronic, mechanical, photocopying, recording or otherwise, without the prior permission of the publisher. All export rights for this book vest exclusively with The Energy and Resources Institute (TERI). Unauthorized export is a violation of terms of sale and is subject to legal action.
Suggested citation Shourie, Abhilasha and Shilpa S. Chapadgaonkar. 2015. Bioanalytical Techniques. New Delhi: TERI
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Preface
Science and technology by far depend upon experimental procedures involving one or more analytical methods. Experiments are conducted not only to predict phenomena, but also to validate the results. Analysis of chemical and biochemical entities is metaphorically the path to achieve the objectives of research in life sciences. In the past few decades, the field of life science has witnessed rapid advancements through development of highly sophisticated, automated, sensitive, and accurate analytical techniques. A vast range of analytical techniques and respective instruments are available, and therefore it is imperative to understand the principles, limitations, and alternatives of a given technique in order to apply it effectively to obtain useful results. ‘Bioanalytical techniques’ is included as a fundamental paper in most courses in chemistry, biochemistry, biology, pharmaceutical and clinical sciences, environmental, forensic and materials sciences. In the past 10 years of our career, we have felt a profound need for a comprehensive yet intensive textbook on analytical techniques that would serve the purpose of both students and researchers. This book has been written to fulfil the need for a text-cumreference in undergraduate and postgraduate level curriculum, providing necessary information required to demonstrate the concepts of an analytical technique in all its guises. The book emphasizes on imparting profound knowledge that is able to meet the current throughput screening demands of scientists and researchers. It consists of 12 chapters, encompassing techniques used for biological and biochemical separation, purification, identification and quantification, all put together to construct a compact package. The chapters have been prepared meticulously using simple yet lucid language. We have included a fairly good
vi Preface number of state-of-the-art sophisticated techniques that would be beneficial to researchers in various fields for experimentation. Many typical analytical procedures appear in the book as boxed features which give a snapshot of the techniques being widely used contemporarily in research. Throughout the book, wherever felt required, extended illustrative examples have been incorporated, which are to be read as part of the respective chapter. The format also encourages the reader to analytical thinking and practical application of the technique. This book will certainly prove to be an invaluable reference tool for students, teachers and researchers in the mentioned fields. We express our gratitude to Almighty and all those who have contributed to the book in any manner. We thank our families for their encouragement and unconditional support. We also thank the editorial team of TERI Press for their continuous effort and faith in us. We solicit constructive suggestions from all the readers for further improvement of the content of the book. Dr. Abhilasha Shourie Dr. Shilpa S. Chapadgaonkar
Contents
Preface
1.
GENERAL PRINCIPLES OF ANALYTICAL INSTRUMENTATION 1.1 1.2 1.3 1.4
2.
Introduction Experimental Studies Experimental Errors Statistical Parameters for Validation of an Experiment
SOLUTIONS AND BUFFERS 2.1 2.2 2.3 2.4 2.5 2.6 2.7
3.
v
Introduction Units of Concentration The Concept of pH Acids and Bases Henderson–Hasselbalch Equation Determination of pKa Buffers
MICROSCOPY 3.1 3.2 3.3 3.4 3.5 3.6 3.7
Historical Background Nature of Light Compound Microscope Image Formation in a Light Microscope Phase Contrast Microscopy Fluorescence Microscopy Electron Microscopy
1 1 2 7 7
11 11 12 16 17 19 20 20
37 37 39 41 43 47 51 54
viii Contents 4.
CELL DISRUPTION 4.1 4.2 4.3 4.4 4.5 4.6 4.7 4.8
5.
CENTRIFUGATION 5.1 5.2 5.3 5.4 5.5 5.6
6.
Introduction Types of Chromatographic Techniques Planar Chromatography Column Chromatography Protein Purification Strategies
ELECTROPHORESIS 7.1 7.2 7.3 7.4 7.5 7.6 7.7 7.8
8.
Introduction Principles of Centrifugation Centrifuge Machines Centrifugal Separations Analytical Ultracentrifugation Care and Safety of Centrifuges
CHROMATOGRAPHIC TECHNIQUES 6.1 6.2 6.3 6.4 6.5
7.
Introduction Barriers for Cell Disruption Methods of Cell Disruption—an Overview Mechanical Methods of Cell Disruption Non-Mechanical Methods of Cell Disruption Combinations of Cell Disruption Methods Selection of Cell Disruption Methods Analysis of Cell Disruption
Introduction Principles of Electrophoresis Free Solution Electrophoresis Paper Electrophoresis Gel Electrophoresis Gel Electrophoresis of Proteins Gel Electrophoresis of Nucleic Acids Capillary Electrophoresis
SPECTROSCOPY I 8.1 Electromagnetic Radiations 8.2 Ultraviolet and Visible Light Spectroscopy 8.3 Fluorescence Spectroscopy
69 69 70 70 71 82 86 88 90
91 91 93 102 109 117 122
125 125 126 127 137 178
181 181 182 187 188 190 199 212 218
225 225 228 238
Contents
8.4 Atomic Absorption Spectrometry 8.5 X-ray Spectroscopy 8.6 Circular Dichroism and Optical Rotatory Dispersion
9.
SPECTROSCOPY II 9.1 9.2 9.3 9.4 9.5
Introduction Infrared Spectroscopy Nuclear Magnetic Resonance Spectroscopy Electron Spin Resonance Spectroscopy Mass Spectroscopy
10. RADIOISOTOPE TECHNIQUES 10.1 10.2 10.3 10.4 10.5 10.6 10.7 10.8
Introduction and Applications Structure of Atom and Radioactivity Types of Radioactive Decay Interaction of Radioactivity with Matter Kinetics of Radioactive Decay Units of Radioactivity Detection and Measurement of Radioactivity Safety Issues and Radio-waste Management
11. IMMUNOCHEMICAL TECHNIQUES 11.1 11.2 11.3 11.4
Introduction Production and Purification of Antibodies Immunoassay Techniques Advances in Immunochemical Techniques
12. MOLECULAR BIOLOGY TECHNIQUES 12.1 12.2 12.3 12.4 12.5 12.6 12.7
Introduction Isolation of Nucleic Acids Qualitative and Quantitative Evaluation of Nucleic Acids Nucleic Acid Hybridization Restriction Digestion Nucleic Acid Sequencing DNA Amplification by Polymerase Chain Reaction
Further Readings Index About the Authors
ix 243 246 253
259 259 259 267 276 282
295 295 297 298 301 303 306 307 320
325 325 329 336 353
355 355 355 364 367 376 380 387 393 397 411
1 General Principles of Analytical Instrumentation
1.1
INTRODUCTION Biological experimentation and chemical experimentation inevitably require analytical instrumentation to perform detection, isolation, purification, identification, and quantification of not only organic and inorganic material, but also of cells and cellular entities. The utility of these techniques and, thereby, instruments span from laboratory scale testing and experimentation to advanced research and commercial applications. Advances in imaging techniques have enabled visualization of micro–nano scale features in biological systems, which have proved beneficial for cellular biology, microbiology, molecular biology, and medical research. Chromatographic techniques enable efficient separation, identification, and measurement of a wide array of compounds such as amino acids, proteins, enzymes, hormones, and drugs on the basis of their biochemical properties. High performance liquid chromatography (HPLC) and gas chromatography (GC) have become indispensable tools in analytical chemistry. The demands for faster and efficient chemical and biochemical separations have led to achieving new horizons in advanced chromatography. Identification and quantification of inorganic and organic compounds, based on their spectral and fluorescent characteristics, can be done with high efficiency and accuracy using spectrometric techniques. Highly sensitive and specific measurement of trace metals and minerals, and structure determination of macromolecules can be done with the help of techniques such as atomic absorption spectrometry, X-ray diffraction, X-ray photoelectron spectroscopy, nuclear magnetic resonance (NMR), and electron spin resonance (ESR). Mass spectrometry (MS) is nowadays an indispensable technique for the molecular analysis of complex compounds. Electrospray-ion trap, matrix-assisted laser
2 Bioanalytical Techniques desorption/ionization time-of-flight (MALDI–TOF), and Fourier transform ion cyclotron resonance (FT–ICR) mass spectrometry are high-resolution sophisticated techniques that enable structural analysis of complex molecules, such as proteins, and reveal their secondary and tertiary conformations. In the past two decades, there have been remarkable improvements in the analytical methods, which have significantly broadened their applications in the analysis of biological and chemical molecules. Hyphenated techniques combining chromatographic and spectral methods have been instrumental in advancements in analytical sciences. Techniques such as liquid chromatography with mass spectrometry (LC–MS) and gas chromatography with mass spectrometry (GC–MS) are being widely used where biochemical separations and subsequent analysis of separated compounds are required. Capillary electrophoresis combined with mass spectrometry (CE–MS) and liquid chromatography combined with nuclear magnetic resonance (LC–NMR) have extended the use of these techniques to structure elucidation and quantification of analyte concentrations as low as nano or pico grams with high resolution and sensitivity. Bioanalytical procedures and instruments are often selective and sensitive in nature and are critical for the successful conduct of the intended experimental analysis. A particular experimental study may require several instruments and methods to be employed in order to obtain the necessary data. Integration of many processes and interdisciplinary methods may need some “standard operating procedures” to define a uniform set of specifications to facilitate interoperability. Often in biological and chemical sciences, uniform nomenclature systems and standardized measurement systems, such as standard international (SI) units, are routinely used. Besides these, it is also a common practice to standardize the experimental conditions, treatments, controls, and other experimental parameters under which the intended data is generated with utmost reliability.
1.2
EXPERIMENTAL STUDIES In science and technology, both elementary studies and research are essentially based on experimental methods. Scientific research is a directional process based largely on keen observations, logical hypothesis formulation, experimental verification of the prediction, and unbiased inferences. Experimental methods involve deliberate manipulation of one or more variables in comparison to other variables or a constant. This approach demands systematic and controlled experimentation involving field, laboratory settings, or both. Experiments generate vital data which after thorough analysis lead to inferences and conclusions that address the scientific query and validate the prediction or hypothesis.
General Principles of Analytical Instrumentation
1.2.1
3
Experimental Design
It is necessary that the experimental procedures and processes be organized properly into a well-formulated “experimental design” to ensure the collection of right type and amount of data depending upon the objectives of the study. A good experimental design, thus, minimizes bias and maximizes reliability of the data so that it tests exactly what the hypothesis states. An experimental unit can be an individual or a group under study. If it is a group, it should be fairly homogeneous because any sort of variations may cause bias and uncertainty in the results. The three basic principles of experimental design: control, randomization, and replication are discussed here at length. • Control Experimental studies are usually carried out under controlled conditions where two or more things are compared to observe the causes or effects of imposed or exposed factors, often termed as “treatment” on the experimental units. The group that is exposed to all the conditions and factors of the experiment except the one which is being tested or is under question is called control group. On the other hand, the group which gets all of the conditions and factors of the experiment along with the variable that is being tested through the experiment is called experimental group. Use of control group eliminates experimental errors and bias. Two types of control groups are recognized: negative control and positive control. In the negative control group, there is no effect of the imposed factors on the subjects, that is, it shows null effect. The positive control group essentially shows an effect of the experimental treatment, conditions, or factors on the subjects being investigated in the experiment. • Randomization It refers to random assignment of experimental units to an experimental group without any prejudice by the experimenter, so that every experimental unit has the same chance of receiving a treatment. If this is done absolutely at random, without imposing any specific selection, then it is called a completely randomized design. If the experimental subjects are first divided into homogeneous blocks which are subjected to the same extraneous variables and then randomly assigned to a treatment group, it is called randomized block design. Randomization is generally used for eliminating bias in experiments. • Replication It is the process of repeating each treatment or measurement on several experimental units (EU) to determine the actual effects of the treatment. Replication helps in identifying the variation among units and eliminating the errors. The statistical inferences obtained by measurements taken from several replicates enhance the reliability and validity of the results.
4 Bioanalytical Techniques 1.2.2
Constants and Variables
The factors or conditions under study in an experiment are categorized as constants and variables. All those factors that remain same for both the experimental group and the control group are called constants and all those factors that vary are called variables. The variable that is varied by the experimenter in a very precise manner is called an independent variable and the one that is measured in response to the independent variable is called a dependent variable or response variable. Each treatment is a combination of one or more independent variables. For example, while studying the effect of fertilizer on crop growth, the amount of fertilizer supplied would be an independent variable, whereas the growth parameters such as height or weight of the plants would be the dependent variables. The controlled variables which would remain same for all the values of independent variable would be the type of plant, the type of fertilizer, temperature, sunlight, humidity, and so on. Variables can be classified as categorical and continuous: • Categorical variables These variables are qualitative and differ in kind, not in magnitude; therefore, they are also called discrete variables. These can be further classified as nominal and ordinal variables. Nominal variables have two or more categories which cannot be quantified and also cannot be placed in an order. For example, gender and marital status. Ordinal variables are variables that have two or more categories which can be ordered or ranked, but the intervals between the scale points may not be even. For example, educational level and socio-economic status of families. • Continuous variables These variables have a numerical value and differ in magnitude, therefore, are also known as quantitative variables. For example, income and age. These can be further classified as interval and ratio variables. An interval variable can be measured along a continuum such as temperature measured in degree Celsius (°C). For example, the temperature of 50°C is higher than 40°C, and an increase from 30°C to 50°C is twice as much as the increase from 40°C to 50°C. Ratio variable is a type of interval variable not measured on a linear scale. In this type, 0 (zero) value of a measurement means absolute absence of such a variable. For example, mass and distance. • Other variables A factor or variable that is not related to the study but may affect a response variable in the study is called an extraneous variable. If the effect of the extraneous variable on the response cannot be distinguished, it is called a confounding variable.
General Principles of Analytical Instrumentation
1.2.3
5
Sampling
The entire population with which the experimental study is concerned, needs to be examined to draw fair conclusions. However, it is often not feasible to measure the effects of treatments on every individual subject from each experimental group, especially when the population under consideration is large. In such a case, the study variables are measured on some randomly selected individuals, called “sample”, which are the subsets of that population. The results from the sample are then used to draw valid inferences about the entire population. The sampling method is thus a scientific procedure of selecting those units from which the required estimates can be obtained. Such results are associated with some degree of uncertainty, as only a part of the population is actually studied and not the whole. Therefore, sampling must always be accompanied by measures of counteracting bias such as random selection of sample and appropriate sample size. The size of the sample depends on the extent of variability in the population. The greater the variability in a population and the larger the sample required, the better it is for estimating characteristics of the population with high degree of accuracy. The basic types of sampling designs are non-probability, probability sampling, systemic sampling, stratified sampling, and cluster sampling. • Non-probability sampling In this type of sampling, units are selected deliberately or on purpose, therefore, it tends to yield biased results. The method gives no assurance of giving equal chance of selection to all the experimental units; hence, it is always associated with high risk of bias. Sampling error cannot be estimated in this sampling. • Probability sampling It is also called “random sampling” in which the selection of sampling units is done without imposing any predictions or preferences. In this method, every individual has an equal probability of selection, therefore, it is free of bias. • Systemic sampling It is done when a starting point is chosen from the sampling frame at random and thereafter every nth member is chosen. This method tends to spread the sample more evenly throughout the population. However, it may possess some hidden bias leading to minor unidentifiable errors. • Stratified sampling It is carried out when the population is not fairly homogeneous and naturally falls into groups such as age groups, income range, etc. In such a case, the representative sample is obtained by first dividing the parent population into strata composed of as far as possible homogeneous units, and then selecting items from each stratum. Stratification makes the estimates of the population parameters more precise and reliable.
6 Bioanalytical Techniques • Cluster sampling It is another type of random sampling in which the entire population is divided into many small sub-populations, each of which is again a cluster of further smaller units. Once the clusters are chosen, simple random sampling can be done. In order to ensure a representative spread across the entire population, each selected cluster should be as dissimilar as possible.
1.2.4
Measurement Scales
In an experiment, numeric data are generated by measuring variables using four different types of scales classified as nominal, ordinal, interval, and ratio. • Nominal scale It is a system of measurement that places the experimental units into mutually exclusive and exhaustive categories, such that every measurement essentially falls into one of the categories. This type of measurement scale is qualitative type and can be assigned an arbitrary value so that each experimental unit can take only a single value out of the given options. For example, nationality, language, and gender. The measure of central tendency for the nominal scale is “mode”. These rankings, however, do not reveal much about the quantitative differences between the subjects. • Ordinal scale An ordinal scale puts the experimental units in a rank order with respect to the variable being assessed, usually representing a hierarchy of levels. However, the ordinal measurements do not have absolute values and the intervals of the scale may or may not be equal in magnitude. Thus, the ordinal scale only measures the qualitative phenomena. For example, ranks of students in a class. The suitable measure of central tendency for the ordinal scale is “median”. • Interval scale An interval scale has equal differences between the scale values, therefore, it provides more quantitative information than the ordinal scale. However, it does not have a true “zero” point, that is it is not able to measure complete absence of the character being measured. For example, the Fahrenheit degree scale used for measuring temperature. Using interval scale, “mean” is taken as the appropriate measure of central tendency and “standard deviation” is commonly used as a measure of dispersion. • Ratio scale It is similar to interval scale as it shows equal differences between the scale values which are quantitatively meaningful. Ratio scale has a true zero point which indicates complete absence of the character. For example, measurement of weight and length.
General Principles of Analytical Instrumentation
1.3
7
EXPERIMENTAL ERRORS Experimental measurements are often susceptible to errors due to many reasons such as defect in an instrument, calibration error in an instrument, procedural error inherent in the adopted method of experimentation, or an unidentifiable error. Since errors lead to deviation in experimental values from the ideal or “true” value, usually comparisons are made between some standard value recognized as “true” value and the experimental value. Therefore, it is important to determine the accuracy of a particular measurement to ascertain the validity of the results. Errors in investigations can be classified as random and systematic errors. • Random errors These are also called precision errors, intrinsic to all experiments, and caused by lack of repeatability in the measurements. These are unpredictable and called indeterminate errors because they cause inconsistency in the measurements, resulting in a scatter in the data about the true value, since there is equal probability of producing measurements that are higher or lower than the “true” value. • Systematic errors These are also called non-random or determinate errors. These are consistent and repeatable errors, caused due to factors that bias the result in one direction such as inherent defect in an instrument, fault in the method or procedure, etc. These errors can be identified and eliminated to some extent or completely either by using an alternative method or by using a standard value of some reference sample.
1.4
STATISTICAL PARAMETERS FOR VALIDATION OF AN EXPERIMENT Analytical procedures and instruments, which are used for intended analytical applications, must meet the characteristic requirements of that application. Moreover, the validation of analysis essentially depends upon the reliability and reproducibility of the experimental procedure and the instruments used for this purpose. In compliance to these analytical procedures and instruments with the intended use, their assessment is done on the basis of certain performance characteristics such as specificity, linearity, range, detection limit, accuracy, and precision. The acceptability of analytical data directly corresponds to the criteria used for validating the method and requires application of standard statistical tests. A measurement system is considered valid if it is both accurate and precise; therefore precision and accuracy are the two fundamental parameters used for such a validation.
8 Bioanalytical Techniques 1.4.1
Precision
The precision of an analytical method refers to the repeatability or reproducibility of measurements of the same quantity on the same sample under the same conditions. When a measurement is repeated several times and a graph is plotted between the number of times of occurrence of a value and the set of values obtained, it is normally bell-shaped with the results scattered symmetrically about a mean value. This type of distribution is called a Gaussian or normal distribution (Figure 1.1). In such a case, the precision of the data set gives an estimate of random error. A measure of variability is “standard deviation” (SD) which shows how closely all the values are clustered around the mean in a set of data. One SD away from the mean in either direction on the horizontal axis (the two shaded areas closest to the centre axis in Figure 1.1) accounts for somewhere around 68 per cent of the values in a data set. Two SDs away from the mean (the four areas closest to the centre areas) account for roughly 95 per cent and three SDs (all the shaded areas) account for about 99 percent values in a data set. Standard deviation (SDx) is computed as _____________ n _ 2
| ÷ sx =
(xi – x) S_____________ i=1 n – 1
|
where n is the number of data points _ x is the mean of xi, and xi is each of the values of the data.
Figure 1.1 A plot showing Gaussian or normal distribution Note On x-axis is the value in question and on y-axis is the number of data points for each value on the x-axis
General Principles of Analytical Instrumentation
9
In the case of systemic errors, the plot is skewed to one side of the mean value, and then increasing the sample size generally increases the precision. The precision is given in terms of deviation from a mean value and is measured as either a “standard error” (SE) of a mean or a “least significant difference” (LSD). In the case of determining LSD, the significance level used should be stated; for example, 5 per cent LSD.
1.4.2 Accuracy The accuracy of an analytical method determines how close is the mean set of test measurements to the standard or true value for that measurement. The “true value” can either be obtained by referring to some standard measurement or by applying population statistics. The population mean is said to be the best estimate of the true value. Accuracy depends on the desired level of confidence in the test and an acceptable “confidence interval” (CI) that relates the sample mean to the population mean. A CI gives a range of values about the sample mean within which there is a given probability that the population mean lies. This probability is determined by the confidence level, which implies that if a large number of random samples are taken from a population and CIs are constructed for each, then 100 per cent of these intervals are expected to contain the population parameter. For example, a “level of confidence” (LOC) of 95 per cent means that, if 100 values are taken from the population, the true mean accuracy of the measurement located within the CI band will be at least 95 out of the 100 tests whereas uncertainty will be associated with the remaining 5 values. The relationship between the two means is expressed in terms of the SD of the data set, the square root of the number of values in the data set and a factor known as Student’s t-test, which can be represented as
| m = x ± ___tsn | _
__
÷
where _ x is the measured sample mean, µ is the population mean, s is the measured SD, n is the number of measurements, and t is the Student’s t-factor.
2 Solutions and Buffers
2.1
INTRODUCTION Most bioanalytical procedures involve preparation of solutions. These solutions constitute the growth media, extraction media, purification media, formulation media, and reagents required for quantification as well as for characterization. Therefore, it is vital for all the concerned people to understand the preparation of solutions and express their concentration in appropriate units. Section 2.2 will familiarize us with different units for expressing concentration of solutions.
2.1.1
Solution
A “solution” is formed when a substance becomes dispersed homogeneously throughout the liquid in molecular form. The substance, called “solute”, is said to dissolve and the liquid is called a “solvent”. Solutions may consist of solids, liquids, or gases dissolved in a solvent. The properties of a solution are uniform throughout the mixture. The presence of a solute in a solvent maybe indistinguishable from the solvent or it may have colour or odour. Table 2.1 summarizes types of solution and their properties.
2.1.2
Properties of a True Solution
A true solution has the following properties: • A true solution is a homogeneous mixture of solute and solvent. • A solute exists as individual molecule or as ion in solution. Typically the diameter of the solute particle is less than one nanometre (10 –9m). • The solute and solvent do not separate on standing or by centrifugation. • True solutions may be colourful but they are always transparent. Light is not scattered by the solution.
12 Bioanalytical Techniques Table 2.1 Types of solutions: true, colloids, and suspensions S. no.
True solution
Colloidal solution
Suspension solution
1.
It is a homogeneous mixture of two or more substances (the solute is extremely small, less than 10 –9 m).
It is a heterogeneous mixture It is a heterogeneous mixture in which one substance is with large particles suspended dispersed in another. The (greater than 10 –7 m). particles are larger than in a true solution (10 –9 –10 –7 m).
2.
Solute particles exist as Particles form groups of ions, The particles form large groups molecules homogeneously atoms, or molecules that are of insoluble particles. dispersed in solvent. evenly dispersed through solvent.
3.
The solution is clear.
4.
The particles do not settle The particles do not settle under The particles settle under gravity under gravity. gravity. on standing.
5.
It cannot be separated using It can be separated by a semi- It can be easily separated by filters. permeable membrane, that is, filtering. cellophane and cell walls.
6.
Examples are air and Examples are mayonnaise, Examples are sand in water, gasoline smog, butter, whipped cream, dust in air, and flour in water and milk
Colloidal solution appears The mixture is cloudy. cloudy.
• Solubility is the ability of the solute to dissolve in a solvent at a particular temperature. It depends on the nature of the solute and solvent, temperature, and pressure. At “saturation point”, the solvent can no longer dissolve any more solute. • Dissolution of an ionic molecule into water results in the formation of an “electrolyte” solution. The ions of the solute will separate in water. These ions are responsible for conduction of electric current in the solution. • The solution has higher osmotic pressure as compared to solvent and it increases with increase in solute concentration. • As the concentration of the solute in the solvent increases, the boiling point of the solution also increases. • The melting point of the solution decreases as the amount of solute is increased in the solution. • A solution of a solid, non-volatile solute in a liquid solvent shows a decrease in vapour pressure above the solution as the amount of solute is increased.
2.2
UNITS OF CONCENTRATION Concentrations of chemicals are routinely expressed in a variety of units. The choice of measurement units to be used in a given situation depends on the
Solutions and Buffers
13
chemical method to be followed in the experiment. It is, therefore, necessary that we become familiar with the units used and methods of converting between different sets of units. To describe quantities that may take on such extreme values, it is useful to have a system of prefixes that accompany the units. Some of the most important prefixes are given in Table 2.2.
2.2.1
Molarity
Molarity (M) is the most common unit for expressing the concentration of a solution in biochemical studies. The molarity of a solution is the number of moles of the solute dissolved per litre of the solution. Simply it is the amount of substance equal to its molecular mass in grams. Molarity of a solution can be calculated as
Molarity =
Mass of solute in g/Gram molecular mass of solute _______________________________________________ Volume of solution in litres
2.1
For example, a solution labelled as 1 M sodium hydroxide (NaOH) has the equivalent of 1 gram mole of NaOH dissolved in 1 litre of solution. Notice that litre is the volume of solution and not of solvent, which is water. 1 M NaOH solution would be made by measuring out the mass of 1 mole of NaOH (40 g) dissolving it in distilled water and making up the volume of the solution to 1 litre.
2.2.2
Molality
Molality (m) is defined as the number of moles of solute per kilogram of solvent. It has its unit as the lower case letter “m” and is read as “molal”. Number of moles of solute Molality = ________________________ Mass of solvent in kg
2.2
It is primarily used when one is dealing with colligative properties of solutions such as freezing point lowering of solvents by solutes. 1 molal solution of sodium chloride (NaCl) is prepared by dissolving 58.5 g (1 mole) of NaCl in 1 kg of distilled water. Table 2.2 Prefixes of commonly used units Prefix
Value –12
Symbol p
Pico
10
Nano
10 –9
n
10
–6
μ
Milli
10
–3
m
Kilo
103
Micro
k
14 Bioanalytical Techniques 2.2.3
Parts per Million
Parts per million (ppm) is a unit of concentration often used when dealing with very small amounts of metal ions and other solutes in water, air, or soil. Grams of solute Parts per million = ________________ × 1,000,000 Grams of solution
2.3
A solution with a concentration of 1 ppm of lead (Pb2+) is equal to 1 mg lead (II) ion per litre of water, or 1 µg of lead (II) ion per millilitre of water (this is because 1 ml water weighs 1 g). A copper salt (Cu2+) solution reported to be 20 ppm would be equivalent to 20 mg of Cu2+ per litre of water or 20 µg per 1 ml. Concentrations can also be expressed by using percentages. Three different types of percentage concentrations are used including mass per cent, volume per cent, and mass/volume per cent.
2.2.4
Mass Per cent
The mass per cent is used for expressing the concentration of a solution when the mass of a solute and the mass of a solution are given. Mass per cent = (Mass of solute)/(Mass of solution) × 100
2.4
2.2.5 Volume Per cent The volume per cent is used for expressing the concentration of a solution when the volume of a solute and the volume of a solution are given. Volume per cent = (Volume of solute)/(Volume of solution) × 100
2.2.6
2.5
Mass/Volume Per cent
The mass/volume per cent is used for expressing the concentration of a solution when the mass of the solute and volume of the solution are given. Volume per cent is often used for expressing the concentration of a liquid solute in a liquid solvent whereas mass per cent is often used for a solid dissolved in a liquid solvent. Mass/Volume per cent = (Mass of solute)/(Volume of solution) × 100
2.2.7
2.6
Mole Fraction
The mole fraction of a substance in a solution is the fraction represented by number of moles of that substance divided by the total number of moles of all substances present in the solution. The sum of mole fractions of each substance present in the solution equals 1.
Solutions and Buffers
15
If the solution is composed of substances A, B, and C, where NA, NB , and NC are the number of moles of A, B, and C, respectively, then X A = Mole fraction of A = (Number of moles of substance A)/(Total number of moles of all substances in solution) NA 2.7 X A = _____________ N A + N B + NC and X A + XB + XC = 1
2.2.8
2.8
Mole Per cent
The mole per cent (of substance A) is mole fraction X A represented in per cent form. Mole per cent (of substance A) = X A × 100%
2.2.9
2.9
Per cent Saturation
Solubility of a solute in any solvent increases with increase in temperature. A solution is said to be unsaturated at a given temperature if it can still dissolve more solute in it at that temperature. A saturated solution is one in which the solute in the solution is in equilibrium with the pure undissolved solute. A supersaturated solution contains more solute in solution than it would ordinarily hold at a given temperature. For example, if in hot water one dissolves all the salt it can possibly solubilize and cools it slowly, the water would contain all the salt that was previously dissolved in it resulting in the formation of supersaturated solution. If a small crystal of salt is introduced into this water, some salt will crystallize out and a saturated solution will remain. The amount of a substance that is dissolved in a solution compared to the amount dissolved in the solution at saturation is expressed as a per cent saturation. Amount of substance that is dissolved Per cent saturation = _______________________________________________ Amount of substance that is dissolved at saturation 2.10
E2.1 EXERCISE Rahul dissolved 10 g of sugar in 250 ml water. Find out the molar concentration of the solution he has prepared? Solution The molecular weight for sugar C12H 22O11 is 342 g in 1 mol. For solving chemical problems, the unit mol/l is the most useful. Thus, the concentration can be calculated as
16 Bioanalytical Techniques 342 g
Æ
1 mol
10 g
Æ
10/342 = 0.0923 moles
Therefore, 250 ml Æ 1 l
Æ
0.02923 moles 29.23/250 moles
Molarity = 0.1169 M The concentration of sugar solution is 0.116 M.
2.3 THE CONCEPT OF pH The dissociation of water into hydroxide and hydrogen ions is a reversible process and can be represented as H2O Æ H+ + OH–
2.11
The equilibrium constant (Keq ) of a chemical reaction is given by the ratio of concentration of products to reactants at equilibrium. For the dissociation of water, we get [H+] [OH– ] Keq = ___________ [H2O]
2.12
The equilibrium constant for the dissociation of water at 25°C has been measured as Keq= 1.8 × 10 –16. This number is small because only a small fraction of water molecules dissociate. The concentration of water can be determined from the fact that 1 mole of water weighs 18 g and 1 litre of water weighs 1000 g. Hence, the concentration of pure water [H 2O] is 1000 g/l/18 g/mol = 55.5 mol/l. By substituting these into Equation (2.12), we get [H+][OH– ] = 1 × 10 –14
2.13
In pure water, the concentrations of hydrogen and hydroxide ions are about the same. Hence by taking the square root of 1 × 10 –14 we find that [H+] and [OH – ] are each about 10 –7 M. This means that 1 litre of pure water contains about one ten-millionth of a mole of hydrogen or hydroxide ions. When substances are dissolved in water, the concentrations of H+ and OH– can change depending upon the ability of the substances to donate or accept hydrogen ions (H+) that is acidity or alkalinity of substances. Acidity and alkalinity are measured with a logarithmic scale called “pH”. pH is defined as
Solutions and Buffers
17
the negative logarithm of the hydrogen ion concentration (–log[H+]). The pH scale ranges from 0 to 14. On the pH scale, the value 7 is considered to be neutral. Substances that can donate hydrogen ions, thus increasing [H+], are acids. Strong acids have pH much lower than 7. Molecules that accept hydrogen ions, thus decreasing [H+], are bases and have pH higher than 7. pH = – log[H+]
2.14
The pH scale allows biologists to define chemical solutions more conveniently by abolishing the need for exponential notation. In the simplest terms, the value of pH simply gives us the value of the exponent of the hydrogen ion concentration. Each one-unit change in the pH on the pH scale corresponds to a ten-fold change in hydrogen ion concentration in the solution.
2.4 ACIDS AND BASES Acids and bases have been defined with different concepts. However, the definition that is better suited to a given situation is usually applied. We will discuss the different concepts describing acids and bases briefly.
Figure 2.1
The pH scale
18 Bioanalytical Techniques 2.4.1 Arrhenius Concept Arrhenius defined acid as any substance that is capable of providing hydrogen ions (or protons) in aqueous solution. A base is defined as a substance containing hydroxyl groups and capable of providing hydroxide ion (OH–) in aqueous solution. According to this definition, a neutralization reaction combines these two ions to form the solvent, that is, water and a salt. This can be illustrated as in Equation 2.15 stating a general reaction between an acid HA and a base BOH. HA + BOH = BA + H2O
2.4.2
2.15
Bronsted–Lowry Concept
According to the definition, an acid is any substance capable of donating a proton (H+) in a solution. A base is any substance that is capable of accepting a proton. A Bronsted acid dissociates into a proton and conjugate base of acid. This can be observed in Equation 2.16 where A– is acting as the conjugate base of HA. HA = H+ + A–
2.16
For example, the ionization of hydrochloric acid (HCl) in water can be viewed as a Bronsted–Lowry acid–base reaction but actually no chemical reaction is taking place. The chemical species formed when a base accepts a proton is called a conjugate acid of the base as illustrated in Equation 2.17 where BH+ is acting as the conjugate base of B. B: + H+ = BH+
2.17
Acids and bases can be classified as strong and weak acids and bases. Strong acids and bases completely dissociate in water. The others are considered to be weak acids and bases. Strong acids HCl is considered as a strong acid because it has a strong tendency to ionize resulting in a large concentration of H+ as HCl completely dissociates into H+ and Cl– contributing to an increase in [H+] in the solution. For example, HCl – hydrochloric acid, HNO3 – nitric acid, H2SO4 – sulfuric acid, HBr – hydrobromic acid, HI – hydroiodic acid, and HClO4 – perchloric acid are strong acids. Strong bases Strong bases like strong acids dissociate completely in water. For example, NaOH – sodium hydroxide, KOH – potassium hydroxide, and Ba(OH)2 – barium hydroxide.
Solutions and Buffers
19
Weak acids and bases ionize only partially in solutions. Examples of weak acids and bases are as follows: Weak acids HCOOH – formic acid, CH3COOH – acetic acid, CCl3COOH – trichloroacetic acid, HF – hydrofluoric acid, HCN – hydrocyanic acid, and H2S – hydrogen sulfide. Weak bases NH3 – ammonia, N(CH3)3 – trimethyl ammonium, C5H5N – pyridine, and NH4OH – ammonium hydroxide.
2.5
HENDERSON–HASSELBALCH EQUATION Dissociation of weak acid Weak acids or bases do not dissociate completely in water and the dissociation is represented by equilibrium. For example, let us study the dissociation of a simple acid (HA). It can be described by the given chemical reaction (Equation 2.18): HA Æ H+ + A–
2.18
The equilibrium constant for the dissociation of a weak acid is given by Equation 2.19: [H+][A– ] Ka = ________ [HA]
2.19
The Henderson–Hasselbalch equation is derived from this by taking the logarithm of both sides and rearranging to get [A– ] – log [H+] = – log Ka + Log _____ [HA]
2.20
We can rewrite Equation 2.20 by replacing –log [H+] with pH and –log Ka with pKa to get Equation 2.21: [A– ] pH = pKa + log _____ [HA]
2.21
We can think of pKa as pH at which the number of molecules of conjugate base [A – ] and weak acid [HA] is equal. The value of log [1] = 0; thus, we get pKa = pH. We can interpret the above equation in the following way: pH is a function depending on the constant (pKa) as well as the ratio of the conjugate base to weak acid ([A– ]/[HA]). Henderson–Hasselbalch equation is
20 Bioanalytical Techniques an extremely useful equation from which pH of the solutions of various ratios of concentrations of conjugate acid and conjugate base form of a substance can be determined. Conversely, it can also be used to find out the required ratio of the conjugate acid-base pair of known pKa to obtain a buffer of desired pH. The concept of buffers is elaborated in the following section.
2.6
DETERMINATION OF pKa Titration leads to the establishment of pKa values. The free acid form of the chemical entity is titrated with an appropriate base. The pH of the solution is monitored as a strong base is added to the buffer solution. This titration curve is recorded as equivalents of base on abscissa against pH as the ordinate. Figure 2.2 shows the titration curve for CH3COOH. The point of inflection indicates the pKa value. Polybasic buffer systems have more than one pKa value. Figure 2.3 shows the titration curve for phosphoric acid, a tribasic acid. The titration of phosphoric acid shows five inflection points where three signify pKa , pKa and pKa , and the other two points indicate where H2PO4 and 1
2
3
HPO4 exist as the sole species.
2.7
BUFFERS Cells must constantly maintain their pH in order to function properly. Enzymes, the universal catalysts of living beings, function at an optimal pH and get inactivated by extremes of pH. In biological systems, control of pH in solutions is a very important part of the practice while working. The control of pH is brought about by buffers. Buffers are pairs of related chemical compounds capable of resisting large changes in pH of a solution caused by addition of small quantities of acids and bases. A buffer system is composed of a weak acid and its conjugate base or a weak base and its conjugate acid. The two components of the system, the buffer pair, complement each other. When small quantities of hydrogen ions are introduced into the medium, they will react with the conjugate base or basic member of the buffer pair to form a weak acid. The weak acid is only slightly ionized. In a similar manner, if small amounts of hydroxide ions are added to the medium, it will react with the weak acid of the buffer pair to form water and conjugate base. In this way large changes in hydrogen ion concentrations are resisted. Let us take the example of acetate buffer. Acetate buffer consists of CH3COOH and CH3COONa. In solution, the dissociation of CH3COOH will be given as
Solutions and Buffers
CH3COOH(aq)
CH3COO –(aq) + H+(aq)
21
2.22
Since CH3COOH is a weak acid, it will not dissociate completely; thus, it has a smaller ratio of ions (CH3COO – and H+) to uncharged acid (CH3COOH) in this equilibrium and a small pKa value. Therefore, the equilibrium is said to be shifted towards the left. The added CH3COONa completely dissociates into ions in the solution. Therefore, according to Le Chatelier’s principle, this will swing the position of the equilibrium even further to the left. Let us understand how buffers resist large changes in pH of the solutions by considering the following situations: Acid is added to a solution containing acetate buffer The buffer solution must remove most of the new hydrogen ions. Otherwise, the pH would drop markedly. Hydrogen ions combine with the acetate ions to make CH3COOH. Although the reaction is reversible, since CH3COOH is a weak acid, most of the new hydrogen ions are removed in this way. CH3COO – (aq) + H+(aq)
CH3COOH(aq)
2.23
Since most of the new hydrogen ions are removed, the pH will not change to a great extent. When alkali is added to this buffer solution Alkaline solutions contain hydroxide ions. There are two mechanisms discussed further in this section by which these hydroxide ions are removed in a buffer solution. (i) Removal of hydroxide ions by reacting with CH 3COOH In the solution containing acetate buffer, CH3COOH molecules are available to react with the added hydroxide ions. They will react to form acetate ions and water. CH3COOH(aq) + DH(aq)
CH3COO –(aq) + H2O (I)
2.24
(ii) Removal of hydroxide ions by reacting with hydrogen ions In a buffer solution there are some hydrogen ions present from the ionization of the CH3COOH. CH3COOH(aq)
– CH3COO(aq) + H+(aq)
2.25
Hydroxide ions can combine with these to make water. The equilibrium shifts more towards right to replace them. This keeps on happening until most of the hydroxide ions are removed.
22 Bioanalytical Techniques
Figure 2.2
Titration curve for CH3COOH
Figure 2.3
Titration curve for phosphoric acid
Solutions and Buffers
23
Equilibrium moves to replace the removed hydrogen ions
CH3COOH (aq)
CH3COOH–(aq) + H+(aq)
Hydroxide ions combine with these to make water
Again, to maintain the equilibrium, not all of the hydroxide ions are removed. Water that is formed re-ionizes to a very small extent to give a few hydrogen ions and hydroxide ions.
E2.2
EXERCISE A solution contains 0.05 M CH3COOH and 0.05 M CH3COONa. Calculate the change in pH when 0.001 mole of HCl is added to a litre of the solution, assuming that the volume increase upon adding HCl is negligible. Compare this to the pH if the same amount of HCl is added to a litre of pure water. Solution Initially, when there is no HCl added to the buffer the hydrogen ions are contributed by the dissociation of CH3COOH. Since CH3COOH is a weak acid, it dissociates only to a small extent to give acetate ions and hydrogen ions. CH3COONa dissociates completely in solution to give acetate ions and sodium ions. Therefore, according to Le Chatelier’s principle the ionization of CH3COOH is further reduced due to common ion effect. The acid equilibrium equation remains the same as follows: Ka =
+ [H ][A– ] ________ [HA]
Since acetate ions contributed by CH3COONa are much higher than the acetate ions contributed by CH3COOH, the concentration of acetate ions [A – ] can be considered to be equal to the concentration of CH3COONa (0.05 M) and the ions contributed by CH3COOH can be neglected. Moreover, since CH3COOH dissociates only to a small extent and its dissociation is further suppressed due to common ion effect, the total concentration of undissociated CH3COOH [HA] can be considered to be constant (0.05 M). Let ‘y’ be the hydrogen ion concentration in the buffer solution. Thus, Ka = y[0.05]/[0.05]
24 Bioanalytical Techniques
y
= Ka
Thus, Ka (CH3COOH) = 1.76 × 10 –5 mole/litre pH = pKa = 4.75 When 0.001 M of HCl is added to the buffer solution, the added protons from HCl combine with the acetate ions to form more CH3COOH, giving the balanced equation of A– + H+ Æ HA [HA] = 0.050 + [H+]HCl = 0.051 mole/litre [A– ]
= 0.050 – [H+]HCl = 0.049 mole/litre
y(0.049) Ka = ________ (0.051) y
= 1.76 × 10 –5 × [(0.051)/(0.049)] = 1.83 × 10 –5 mole/litre
pH = 5 – 0.26 = 4.74 In the presence of a buffer, the pH changes from 4.75 to 4.74, a difference of only 0.01 units. In the absence of the buffer, if 0.001 moles of HCl is added to 1 litre of water, the pH would be pH of 3. Therefore, it can be said that the buffer resists large changes in pH.
2.7.1
Preparation of Buffers
Method for preparation of buffers is available as buffer tables in numerous books. Also many websites have “buffer calculators” that help us to calculate the quantities of buffering species required to prepare buffers of desired pH and concentration. Henderson–Hasselbalch equation is the basis of computation in such tables. The method is straightforward and has been exemplified as follows: Let us prepare 250 ml of 0.05 M phosphate buffer at pH 7.5. Remember Henderson–Hasselbalch equation (Equation 2.20) [A– ] _____ pH = pKa + log [HA]
Solutions and Buffers
25
where A– and HA are the concentrations of basic and acidic components of the buffer system, respectively. So, let us calculate the total ion concentrations of both the components of the buffer system, that is for phosphate buffer, the components are monobasic potassium phosphate and dibasic potassium phosphate. For 1000 ml we require 0.05 moles of total ion concentrations in both the solutions. Therefore, for 250 ml we would require 250 ml × 0.05 M/1000 ml = 0.0125 M or 12.5 mM in both the solutions. According to Henderson–Hasselbalch equation, for the phosphate buffer system of pH 7.5 concentration ratio is given by 7.50 = 6.86 + log ([K 2HPO4]/[KH2PO4]) Therefore, [K 2HPO4]/[KH2PO4] = 4.37 Let the concentration of KH 2PO4 in the buffer be ‘a’ molar and the concentration of K 2HPO4 be ‘b’ molar which is equal to 4.37a. However, the total ion concentration has to be 0.125 M. Therefore, 0.125 M = 4.37a + a Converting to target buffer solution concentration and units, a = 0.0233 moles of KH2PO4 = 0.0233 mol × 136.086 g/mol = 3.17 g b = 0.1016 moles of K 2HPO4 = 0.1016 mol × 174.176 g/mol = 17.70 g Thus, we need to weigh 3.17 g of monobasic phosphate and 17.6904 g of dibasic phosphate in 250 ml of deionized water. Note: It would be advisable to mix the two solutions in little less than the desired volume and check the pH. If it is somewhat different, then the pH may be adjusted by adding acid or base making the buffer of exact pH.
2.7.1.1
Buffer concentration
Buffers are effective when used at adequate concentration. Let us understand this concept by considering the following example: Consider two different buffer solutions A and B. Buffer A contains 0.01 M CH3COOH and 0.01 M CH3COONa whereas buffer B contains 0.1 M CH3COOH and 0.1 M CH3COONa. These two buffer solutions were titrated with HCl and the change in pH was calculated for a given amount of HCl
26 Bioanalytical Techniques (assuming that the volume increase upon adding the HCl is negligible). The results are tabulated in Table 2.3. Table 2.3 Effect buffer concentration on buffer capacity Moles of HCl added
pH of buffer A [0.01 M CH3COOH and 0.01 M CH3COONa]
pH of buffer B [0.1 M CH3COOH and 0.1 M CH3COONa]
0.000
4.75
4.75
0.001
4.67
4.75
0.002
4.58
4.74
0.003
4.49
4.73
0.004
4.39
4.72
0.005
4.28
4.71
The pH of buffer A changes more rapidly as compared to the pH of buffer solution B. Therefore, it can be said that increasing the buffer concentration improves the ability of the buffer to resist the changes in pH. An adequate buffer capacity is often reached at concentrations higher than 25 mM. Optimum concentrations lie between 20 mM and 100 mM. However, ionic strength may affect enzyme activity, therefore, caution has to be exercised.
2.7.1.2
Buffer capacity
The ability of a buffer system to resist pH changes is its buffer capacity and indicated by the buffer index (b): b = DB/DpH
2.26
where B is the number of moles of strong base added and DpH is the change in pH. Buffer capacity is defined as the number of equivalents of strong base (DB) required for a unit change in pH (DpH) in 1 litre of solution. With the increase in the buffer capacity, the change in pH decreases when a given amount of strong acid or base is added to it. Buffer capacity is dependent on the total concentration of the buffer system and on the HA/A – ratio. The buffer index number is generally experimentally derived in a manner like titration. Buffer capacity is maximum when pH = pKa and is acceptable in the range pH = pKa ± 1.
Solutions and Buffers
2.7.1.3
27
Effect of temperature on pH
Generally when we consider the use of buffers, we make the following two assumptions: 1. The activity coefficient of the buffer ions is approximately equal to one over the useful range of buffer concentrations. 2. The value of Ka is constant over the working range of temperature. However, in real practice one observes that pH changes slightly with change in temperature. This might be very critical in biological systems where a precise hydrogen ion concentration is required for reaction systems to operate with maximum efficiency. The difference might appear to be slight but it has significant biological importance. Although the mathematical relationship of activity and temperature may be complicated, the actual change of pKa with temperature (DpKa /°C) is approximately linear.
2.7.1.4
Effects of buffers on factors other than pH
Buffer ions may interfere with the other components under study. For example, citric acid and citrates are potential calcium chelators; therefore, citrate buffers are regarded as unsuitable for studies involving calcium ions. If the system under study contains phosphates, the phosphates react with calcium resulting in production of insoluble calcium phosphate precipitates. Phosphate ions in buffers may inhibit the activity of enzymes, such as carboxypeptidase, fumarease, carboxylase, and phosphoglucomutase. Tris (hydroxy–methyl) aminomethane can chelate copper and also acts as a competitive inhibitor of some enzymes. Tris-based buffers are not recommended when studying the metabolic effects of insulin. Moreover, buffers like Tris that contain primary amine groups may interfere with the Bradford dye-binding method of protein assay. Buffers such as ACES, BES, and TES have a tendency to bind copper. HEPES and HEPPS buffers are not preferable when a protein assay is performed by using Folin reagent. Borate buffers are not suitable for gel electrophoresis of protein; they can cause spreading of the zones if polyols are present in the medium.
2.7.1.5
Contamination in buffers
Contamination in buffers can be prevented by (i) sterilization by filtration or autoclaving and storage in sterile conditions, (ii) addition of 0.02% w/v sodium azide, (iii) storage at +4°C, and (iv) use of highly concentrated stock solution.
2.7.2
Criteria for Selecting Buffers (i) Buffer should possess sufficient buffering capacity in the desired pH range Generally, buffers are most effective over a range of one
28 Bioanalytical Techniques pH unit on either side of their pKa value. For example, Tris that has a pKa value of 8.3 has an effective pH range of 7–9. If you expect the pH to drop during the experiment, choose a buffer with a pKa slightly lower than the working pH. This will permit the buffering action to become more resistant to changes in [H+] as hydrogen ions are liberated. Conversely, if you expect the pH to rise during the experiment, choose a buffer with a pKa slightly higher than the working pH. For best results, the pK a of the buffer should not be affected significantly by buffer concentration, temperature, and the ionic constitution of the medium. The pKa values of some important buffers have been given in Table 2.4. (ii) Minimum effect of temperature and ionic concentration on buffer pH The buffer pH should not vary considerably on change in temperature or ionic concentration during the experiment. (iii) Inert It should be chemically inert and should not react or bind with biomolecules. Tris is known to be a biological inhibitor and reacts with primary amines. (iv) Purity It should be available in high degree of purity and should not contain impurities that interfere with the analysis. “Buffers” prepared from reagent grade chemicals may contain UV absorbing “impurities” that interfere with UV liquid chomatographic detectors. (v) Stable It should not undergo enzymatic or hydrolytic degradation. For example some buffers are not stable with Osmium tetraoxide (OsO4) or aldehydes at high concentrations. (vi) Non-toxic Buffer should be tissue compatible if used as a rinse or to store the tissue for any length of time. Otherwise, toxic effects will be seen (dead or stressed tissue or autolysis before fixation). Tris can penetrate biological membranes in its unionized form. Tris buffer reacts with primary amines, and modifies electron transport and phosphorylation in chloroplasts. Tris also inhibits respiratory enzymes in mitochondria. Typical “vital buffers” include PBS, Tris, Hepes, and Pipes. (vii) No absortion of light This is very important as absorption of light would interfere with spectrophotometric analysis.
2.7.3
Commonly Used Buffers
The pKa values of commonly used buffers are given in Table 2.4. The useful pH range is within pKa ± 1 of the buffer.
Solutions and Buffers
2.7.4
29
Recipe of Commonly Used Buffers
The information provided below is intended only as a general guideline. We strongly recommend the use of a sensitive pH meter with appropriate temperature setting for final pH adjustment. Addition of other chemicals, after adjusting the pH, may change the final pH value to some extent. The buffer concentrations in the tables below are used only as examples. You may select higher or lower concentrations depending upon your experimental need. 1. Hydrochloric acid–Potassium chloride buffer (HCl–KCl); pH range 1.0–2.2
Table 2.4 pKa values of commonly used buffers Buffer
pK a
Buffer
pK a
H3PO4 / NaH2PO4 (pK a1)
2.12
Bicine
8.35
Glycine (pK a1)
2.34
Glycylglycine (pK a2)
8.40
Citric acid (pK a1)
3.13
TAPS
8.40
CH3COOH
4.75
Bis-Tris Propane (pK a2)
9.00
Citric acid (pK a2)
4.76
Boric acid (H3BO3 / Na2B4O7)
9.24
MES
6.15
CHES
9.50
Cacodylic acid
6.27
Glycine (pK a2)
9.60
H2CO3 / NaHCO3 (pK a1)
6.37
NaHCO3 / Na2CO3 (pK a2)
10.25
Citric acid (pK a3)
6.4
CAPS
10.40
Bis-Tris
6.50
Piperidine
11.12
ADA
6.60
Na2HPO4/Na3PO4 (pK a4)
12.67
Bis-Tris Propane (pK a1)
6.80
PIPES
6.80
ACES
6.90
Imidazole
7.00
BES
7.15
MOPS
7.20
NaH2PO4 / Na2HPO4 (pK a2)
7.21
TES
7.50
HEPES
7.55
HEPPSO
7.80
Triethanolamine
7.80
Tricine
8.10
Tris
8.10
Glycine amide
8.20
30 Bioanalytical Techniques (a) 0.1 M KCl: 7.45 g/l (MW: 74.5) (b) 0.1 M HCl Mix 50 ml of KCl and indicated volume of HCl. Mix and adjust the final volume to 100 ml with deionized water. Adjust the final pH using a pH meter. ml of HCl pH
97
64.5
41.5
26.3
16.6
10.6
6.7
1.2
1.4
1.6
1.8
2.0
2.2
1.0
2. Glycine–HCl buffer; pH range 2.2–3.6 (a) 0.1 M Glycine: 7.5 g/l (MW: 75) (b) 0.1 M HCl Mix 50 ml of glycine and indicated volume of HCl. Adjust the final volume to 100 ml with deionized water. Adjust the final pH to desired pH using pH meter. ml of HCl pH
44.0
32.4
24.2
16.8
11.4
8.2
6.4
5.0
2.2
2.4
2.6
2.8
3.0
3.2
3.4
3.6
Citrate buffer; pH range 3.0–6.2 (a) 0.1 M citric acid: 19.21 g/l (MW: 192.1) (b) 0.1 M sodium citrate dihydrate: 29.4 g/l (MW: 294) Mix citric acid and sodium citrate solutions in the proportions indicated and adjust the final volume to 100 ml with deionized water. Adjust the final pH to desired value. The use of pentahydrate salt of sodium citrate is not recommended. ml of citric acid
46.5
40.0
35.0
31.5
25.5
20.5
16.0
11.8
7.2
ml of sodium citrate
3.5
10.0
15.0
18.5
24.5
29.5
34.0
38.2
42.8
pH
3.0
3.4
3.8
4.2
4.6
5.0
5.4
5.8
6.2
4. Acetate buffer; pH range 3.6–5.6 (a) 0.1 M CH3COOH (5.8 ml made to 1000 ml) (b) 0.1 M CH3COONa; 8.2 g/l (anhydrous; MW: 82) or 13.6 g/l (trihydrate; MW: 136) Mix CH3COOH and CH3COONa solutions in the proportions indicated and adjust the final volume to 100 ml with deionized water.
31
Solutions and Buffers ml of CH3COOH
46.3
41.0
30.5
20.0
14.8
10.5
4.8
ml of CH3COONa pH
3.7
9.0
19.5
30.0
35.2
39.5
45.2
3.6
4.0
4.4
4.8
5.0
5.2
5.6
Citrate–Phosphate buffer; pH range 2.6–7.0 (a) 0.1 M citric acid; 19.21 g/l (MW: 192.1) (b) 0.2 M dibasic sodium phosphate; 35.6 g/l (dihydrate; MW: 178) or 53.6 g/l (heptahydrate; MW: 268) Mix citric acid and sodium phosphate solutions in the proportions indicated and adjust the final volume to 100 ml with deionized water. Adjust the final pH using a sensitive pH meter. ml of citric acid
44.6
39.8
35.9
32.3
29.4
26.7
24.3
22.2
19.7
16.9
13.6
6.5
ml of sodium phosphate
5.4
10.2
14.1
17.7
20.6
23.3
25.7
27.8
30.3
33.1
36.4
43.6
pH
2.6
3.0
3.4
3.8
4.2
4.6
5.0
5.4
5.8
6.2
6.6
7.0
Phosphate buffer; pH range 5.8–8.0 (a) 0.1 M sodium phosphate monobasic; 13.8 g/l (monohydrate, MW: 138) (b) 0.1 M sodium phosphate dibasic; 26.8 g/l (heptahydrate, MW: 268) Mix sodium phosphate monobasic and dibasic solutions in the proportions indicated and adjust the final volume to 200 ml with deionized water. The final pH is adjusted to the desired pH using pH meter. ml of sodium phosphate 92.0 monobasic
81.5
73.5
62.5
51.0
39.0
28.0
19.0
13.0
8.5
5.3
ml of sodium phosphate dibasic
8.0
18.5
26.5
37.5
49.0
61.0
72.0
81.0
87.0
91.5
94.7
pH
5.8
6.2
6.4
6.6
6.8
7.0
7.2
7.4
7.6
7.8
8.0
Tris–HCl buffer, pH range 7.2–9.0 (a) 0.1 M Tris (hydroxymethyl) amino methane; 12.1 g/l (MW: 121) (b) 0.1 M HCl Mix 50 ml of Tris (hydroxymethyl) amino methane and indicated volume of HCl and adjust the final volume to 200 ml with deionized water. Adjust the final pH using a pH meter.
32 Bioanalytical Techniques ml of HCl pH
44.2
414
38.4
32.5
21.9
12.2
5.0
7.2
7.4
7.6
7.8
8.2
8.6
9.0
8. Glycine–Sodium hydroxide, pH range 8.6–10.6 (a) 0.1 M glycine; 7.5 g/l (MW: 75) (b) 0.1 M sodium hydroxide; 4 g/l (MW: 40) Mix 50 ml of glycine and indicated volume of sodium hydroxide solutions, and adjust the final volume to 200 ml with deionized water. Adjust the final pH. ml of sodium hydroxide
4.0
8.8
16.8
27.2
32.0
38.6
45.5
pH
8.6
9.0
9.4
9.8
10.0
10.4
10.6
9. Carbonate–Bicarbonate buffer, pH range 9.2–10.6 Special note for buffers containing sodium hydrogen carbonate (sodium bicarbonate): this buffer substance requires a closed system. In aqueous solutions, sodium hydrogen carbonate degrades into carbon dioxide (CO2) and sodium carbonate above 20°C. Complete degradation occurs at 100°C. Solutions containing sodium hydrogen carbonate cannot therefore be autoclaved, but have to be sterile filtered. When preparing, they should not be stirred too vigorously and too long. The pH of a freshly prepared 100 mM solution is 8.3 at 25°C. (a) 0.1 M sodium carbonate (anhydrous), 10.6 g/l (MW: 106) (b) 0.1 M sodium bicarbonate, 8.4 g/l (MW: 84) Mix sodium carbonate and sodium bicarbonate solutions in the proportions indicated and adjust the final volume to 200 ml with deionized water. Adjust the final pH. ml of sodium carbonate ml of sodium bicarbonate pH
4.0
9.5
16.0
22.0
27.5
33.0
38.5
42.6
46.0
40.5
34.0
28.0
22.5
17.0
11.5
7.5
9.2
9.4
9.6
9.8
10.0
10.2
10.4
10.8
Solutions and Buffers
33
GOOD’S BUFFERS Good buffers (Good’s buffers) are 12 buffering agents selected and described by Norman E. Good and colleagues in 1966. The properties that make them useful candidates for biochemical research are listed as follows: • pKa values between 6–8 Since most biological reactions take place at nearneutral pH. • Good aqueous solubility and poor non-polar solubility For ease of handling aqueous biological systems, and poor non-polar solvent solubility to prevent the buffer from accumulating in cell membranes and other non-polar compartments of the biological system. • Membrane impermeability Buffers do not readily pass through cell membranes and show a selective permeability reducing the accumulation of the buffer compound within cells. • Minimal salt effects Highly ionic buffers may cause problems or complications in some biological systems. • Well-behaved cation interactions If the buffers form complexes with cationic ligands, the complexes formed should remain soluble. Ideally, at least some of the buffering compounds will not form complexes. • Stability The buffers should be chemically stable, resisting enzymatic and non-enzymatic degradation. • Optical absorbance Buffers should not absorb the visible or ultraviolet spectrum so that it does not interfere with commonly used spectrophotometric assays.
PHYSIOLOGICAL BUFFERS All the physiological processes have to be carried out at a specific pH range or the metabolic processes would be affected. The pH of human blood is closely maintained at 7.4 and if the blood pH falls below 7 or rises to 7.8, it may result in death within minutes. pH plays an important role in almost all biological processes. Small change in pH that is a decrease or an increase can cause metabolic acidosis or alkalosis. Metabolism is usually associated with the release of protons (H+) (decrease in pH) or uptake of protons (H+) (increase in pH). Presence of physiological buffers ensures the maintenance of physiological pH. Important buffers that are dominant in human body are: (i) Bicarbonate buffers, (ii) Phosphate buffers, and (iii) Protein buffers. Contd...
34 Bioanalytical Techniques
E2.3 EXERCISE 1. Determine the molarity of a solution made by dissolving 20 g of NaOH in sufficient water to yield a 482 cm3 solution. 2. If the 0.131 M sugar solution has a density of 1.10 g/ml, what is the concentration in mole percentage? 3. How much NaOH is contained in 25 ml of a solution whose concentration is 0.1234 M? 4. Neeta dissolved 13 g of NaCl in water to make 2 litre of aqueous solution. Calculate the molarity. 5. A solid mixture contains 300 kg NaCl and 400 kg KCl. Find the composition of mixture in (i) mass per cent and (ii) mole per cent. 6. The amount of oxygen dissolved in 1000 g fermentation medium was determined to be 0.007 g. Calculate the dissolved oxygen concentration of medium in ppm.
Solutions and Buffers
7. 8. 9. 10. 11. 12. 13.
14.
15.
16. 17. 18. 19. 20.
21. 22.
35
Calculate the total mass of solute in 1 kg of 5 ppm solution. What is the pH of a 0.3 M solution of CH3COOH? (Ka = 1.8 × 10 –5) What is the pH of a 0.02 M HCl solution? Calculate the pH of a 0.03 M solution of an acidic drug at 25°C (pKa = 4.76 at 25°C). A soap solution has a pH of 11.73. What is the [OH– ]? What is the pH of human intramuscular fluid in which the hydroxide ion concentration is 6.18 × 10 –8 M? Given the need to prepare a solution with its pH buffered to 7.4, which of the following monoprotic acids and bases would be the best to use based on their pKa ? (Methylamine pKa = 10.62, Formic acid pKa = 3.6, Tris pKa = 8.3, TAPS pKa = 8.5, MES pKa = 6.2, TES pKa = 7.6) Consider two buffered solutions P and Q, both buffered with same buffering agent to pH 7.4. The concentration of buffering agent in solution P is 10 mM while in solution B it is 100 mM. If you add 1 ml of 1 N HCl to 100 ml of each of solution P and solution Q, which solution will have the lower pH? Why? In a pH 8 solution of Ethanolamine (HOCH2CH2NH2), which has a pKa of 9.5, which form is in excess: the neutral form or the cationic form (or neither)? Why? What form of Pyruvic acid (anionic or neutral) is predominant in a solution of pH 4.5? (pKa of 4.5) Why? Consider a buffer made from 0.1 M Tris (pKa of 8.3) adjusted to 7.3. What is the ratio of conjugate acid of Tris to the base form of Tris? Calculate the pH of a buffer solution made from 0.20 M CH3COOH and 0.50 M CH3COONa. (Ka CH3COOH = 1.8 × 10 –5) What ratio of acetate ion to CH3COOH would you need to make up an acetate buffer with a pH of 5.25? (pKa = 4.74) What is the concentration of CH3COONa to be added to 0.12 moles of CH3COOH to make a litre if we want a buffer with a pH of 4.50? (pKa = 4.74) What is the pH of a buffer containing equimolar concentrations of ammonia and ammonium chloride? What is its pH? (Kb = 1.77 × 10 –5) An enzyme-catalysed reaction is performed in 100 µl of 0.45 M Tris buffer. Tris exists in a deprotonated (Tris) and a protonated form (Tris–HCl). The pH of the mixture at the start of the assay is 7.20 and the pKa of Tris is 8.10. During the reaction 5 µmoles of H+ ions are consumed. What is the ratio of Tris/Tris–HCl at the start of the reaction? What are the concentrations of Tris and Tris–HCl at the end of the reaction?
3 Microscopy
3.1
HISTORICAL BACKGROUND The science of lenses facilitated the desire to see beyond the human power of vision and led to the development of the art of “seeing the unseen”. Microscopy eventually became an instrumental requirement in many magnificent discoveries in science such as “cytology” (study of cell). In 1665, Robert Hooke used a compound microscope consisting of two lenses to visualize “cells” in a thin section of a cork (Figure 3.1a). Contemporarily, Antonie van Leeuwenhoek discovered bacteria, yeast, and numerous other micron-scale biological entities with the help of lenses. He later on worked on polishing the curvatures of lenses for obtaining better magnifications and resolutions, and built numerous microscopes for studying various biological specimens (Figure 3.1b). Many biologists of early times such as Louis Pasteur and Robert Koch were benefited by the magnifying properties of lenses as it led to many magnificent discoveries in biological science. These applications of lenses laid the fundamental stone in the development of compound microscope (Figure 3.1c). It consists of two convex lenses aligned serially in a tube—an objective lens that lies closer to the object and an eyepiece or ocular lens that lies closer to the eye of the observer. Commercial manufacturing of microscopes started around 1847, when Carl Zeiss introduced simple microscopes consisting of only one lens and thereafter launched microscopes with two lenses. Compound microscopes developed at that time suffered from the defect of chromatic and spherical aberrations due to the use of two or more lenses. The images produced by these instruments were often blurred and had colourful halos which interfered with resolution.
38 Bioanalytical Techniques
Figure 3.1
Outline sketches of some early microscopes: (a) Hooke microscope ca. 1670, (b) Leeuwenhoek microscope ca.1600s, (c) Janssen compound microscope ca. 1600s
High quality objectives with corrected optical aberrations were developed after Ernst Abbe postulated the “Abbe sine condition” in 1872, and for the first time defined the term “numerical aperture”. d = l/2NA 3.1 where d is the resolution or the smallest resolvable distance between two point objects, l is the wavelength of light, and NA is the numerical aperture. This discovery revolutionized optical research by mathematically defining the properties of lenses and proved to be indispensable in designing the objectives of desired specifications.
Microscopy
39
Another milestone in the history of microscopy was the discovery of an illumination method contributed by August Kohler in 1893 which allowed microscopists to use the resolving power of Abbe’s objectives to the fullest. Three subsequent decades were witness to many innovations in light microscopy. Zeiss continued to revolutionize microscope design and in 1933 offered the standing model of microscope with curved arm, horizontal stage, inclined viewing head, and convenient adjustment settings. This model is very popular even today due to its ease of handling.
3.2
NATURE OF LIGHT Visible light represents about 2.5 per cent part of the spectrum of electromagnetic radiation within the wavelength range 400–800 nm, represented by seven colours commonly termed as VIBGYOR. The wave nature of light is characterized by its wavelength or frequency and amplitude (Figure 3.2). Wavelength is defined as the distance between two consecutive crests or troughs of a sinusoidal wave. Frequency measures the number of wavelengths passed through a fixed point in a given time. The energy (E) possessed by the light wave measured in electron volt (eV) is given by the relation: hc 3.2 E = hv = __ l where h is the universal Planck’s constant = 4.135 × 10 –15 eV and c is the speed of light in vacuum = 3 × 108 m/s. The properties of light by virtue of its wave nature are evident on interaction with matter, and some of which are discussed here: (i) Refraction The speed of light in vacuum is constant. When light wave enters any other medium such as air or glass, the speed of its propagation is slowed down and the direction of its propagation is also changed at the boundary of two media having different densities. If c0 is the speed of light in vacuum and c in a given medium, then the index of refraction
Figure 3.2
Diagrammatic representation of a light wave
40 Bioanalytical Techniques or refractive index (n) of that medium can be calculated by the following relation: 3.3 n = c0 /c If the light enters a medium having a higher index of refraction or refractive index, the speed of propagation is slowed down and the refracted light wave bends towards the normal (perpendicular to the surface of incidence). The refraction of light depends upon the refractive index of the medium, the wavelength of light, and the angle of incidence of the light. (ii) Dispersion When white light passes through some material such as prism or lens, it separates spatially into its component wavelengths. Glass prisms, thus, generate different colours of light and, therefore, are used to construct spectrometers. However, while using lenses, this phenomenon generates chromatic aberrations in optical devices due to failure of different wavelengths to focus on the same point. (iii) Diffraction It is the phenomenon of spreading of light around edges, when it passes through an object or aperture (hole or slit). Diffraction is observable when the size of an object or aperture is smaller than the wavelength of the impinging wave. On passing through a slit, light diffracts into a series of circular waves and the wavefront which emerges from the slit is a cylindrical wave of uniform intensity. The angle of diffraction at which new wavefronts emerge from the aperture is given by the following relation: q = l/D
3.4
where q is the angle of diffraction, l is the wavelength of light, and D is the diameter of the aperture or slit. (iv) Interference Interference is observed when two light waves emitted from a coherent source superimpose to form a resultant wave of greater or lower amplitude. Interference is closely associated with diffraction and is observable during diffraction from two slits, each slit acting as a coherent source of light. The diffracted beams from the two slits overlap, causing the waves to superimpose. The interference pattern due to the superposition of the waves is seen as alternate dark and bright bands on the screen, called fringes. If the crests and troughs of the two waves superimpose and the resultant wave is of greater amplitude, it is said to be a constructive interference and is seen as a bright fringe (Figure 3.3a). On the other hand,
Microscopy
Figure 3.3
41
(a) Constructive interference and (b) destructive interference
if a crest of one wave meets a trough of another wave, then the sum of the amplitudes of the two waves will result into a wave with amplitude lesser than either of the waves or even zero. This is called destructive interference and forms dark fringes on screen (Figure 3.3b).
3.3
COMPOUND MICROSCOPE The essential features of a compound microscope as illustrated in Figure 3.4 are as follows: (i) Base The microscope stands firmly on the base and houses an illuminator or a light source. It also consists of a power switch to control the illuminator. (ii) Body tube The body tube of the microscope is attached at the top of the arm and is positioned over the microscope’s stage. It serves to align the eyepiece to the objective lenses. (iii) Arm The arm is meant to connect the body tube to the base of the microscope. It is curved so as to position the body of the microscope over the stage. It consists of focusing knobs near the base. (iv) Illumination Usually light source in a microscope is an electric bulb. Earlier mirrors were used to reflect sunlight to view the specimen. (v) Condenser It consists of a lens system that gathers light from the illuminator and focuses it onto the specimen being viewed. (vi) Iris diaphragm It is located in the condenser and serves to regulate the amount of light reaching the specimen so that it gets adequate illumination.
42 Bioanalytical Techniques
Figure 3.4
A compound microscope
(vii) Aperture It is a hole in the middle of the stage that allows light from the source to reach the specimen. (viii) Stage It is a mechanically controlled platform that holds the specimen firmly through metal clips. It consists of knobs that allow its restricted movements so that all the areas of the specimen can be viewed by little displacement. (ix) Specimen The specimen is usually mounted on a glass slide, fixed either temporarily or permanently using a cover slip on the slide. (x) Nosepiece It is located at the bottom end of the microscope’s body and holds multiple objective lenses. It is a rotating structure that helps to select an objective lens of desired magnification and aligns it to the eyepiece. (xi) Objective lenses An objective lens lies closest to the specimen and is used for its primary magnification on proper focusing. A standard microscope has three to four objective lenses ranging from 4X to 100X
Microscopy
43
of magnification power. The working distance between the objective and the specimen decreases with increasing magnification of objective and is adjusted with the help of focusing knobs. (xii) Eyepiece The eyepiece or ocular lens is used to visualize the magnified image of the specimen. It is usually available with 10X or 15X of magnification power. There is one eyepiece in a monocular microscope and two in a binocular microscope for facilitating binocular vision. Trinocular microscopes are also available having three oculars meant for photographic imaging along with direct viewing of the specimen. (xiii) Focusing knobs These are attached to the sides of the microscope arm. The coarse adjustment knobs are bigger and used for bringing the specimen into general focus, whereas the fine adjustment knobs are smaller and used for fine tuning the focus to visualize more details of the specimen.
3.4
IMAGE FORMATION IN A LIGHT MICROSCOPE Light passing through the specimen is gathered by the objective which projects a real, inverted, and magnified image. This image is focused at a fixed plane within the microscope termed as the intermediate image plane located within a fixed diaphragm in the eyepiece. The distance between the back focal plane of the objective and the intermediate image is called the optical tube length, whereas the distance between the nosepiece and the eyepieces (oculars) in the observation tubes is known as the mechanical tube length of a microscope. The eyepiece further magnifies the real image projected by the objective to form a highly magnified virtual image. A compound microscope consists of two convex lenses, an objective and eyepiece. The objective is placed near the object and has very small focal length fo and small aperture, whereas the eyepiece is close to the eyes of the viewer and has relatively large focal length fe and large aperture. When a small object (O) is placed at distance (x) to the objective, a real, inverted, and magnified image (I) is formed at distance (y). This image lies essentially within the focus of the eyepiece and is used as an object to form a virtual and much magnified final image (I¢) (Figure 3.5). There are three essential requirements for production of a good image in a compound microscope: • Magnification • Resolution • Contrast
44 Bioanalytical Techniques
Figure 3.5
Ray diagram illustrating image formation in a compound microscope
Figure 3.6
Ray diagram illustrating simple magnification by a lens
(i) Magnification In a conventional compound microscope, two converging lenses are used for a two-step magnification process, an objective lens and an eyepiece or ocular lens. The objective projects a real, inverted, and magnified image which is used as an object by the eyepiece to form a highly magnified final image (Figure 3.6). The linear magnification is given by the ratio of the image size to the object size or image distance to the object distance. h¢ –I M = __ = __ O h where M is the linear magnification,
3.5
Microscopy
45
I is the distance of the image from the centre of curvature, O is the distance of the object from the centre of curvature, h¢ is the height of the image, and h is the height of the object. Lens equation expresses the relationship between object distance, image distance, and focal length and is represented as follows: 1 1 1 __ __ __ 3.6 o + i = f where o is the object distance from the centre of curvature, i is the image distance from the centre of curvature, and f is the focal length. The final magnification through the microscope is dependent on the magnifying power of both the lenses; therefore, total magnification (Mtotal) of a microscope can be calculated as the product of magnification of all the lenses (Equation 3.7). Mtotal = Mobjective × Mocular
3.7
Microscope lenses are commonly available with magnification power ranging from 4X to 100X for objectives and 10X to 15X for oculars. (ii) Resolution The light microscope produces images of the specimen under observation by transmitting visible light through it. The phenomena of diffraction, scattering, and interference of light waves largely hamper the resolution. The diffraction of light by the specimen limits the resolution of lenses because the light is not focused to a point, but instead forms an “airy disk” (Figure 3.7). The airy disk has a central bright spot with concentric circles of light of reducing brightness towards the periphery, which makes it difficult to distinguish the two point objects. The larger the aperture of the lens
Figure 3.7
Airy disk hampering resolution of two point objects
46 Bioanalytical Techniques and the smaller the wavelength of light, the finer is the resolution of a microscope. The smallest resolvable distance (d) between the two points is given by Rayleigh criterion. d = 1.22 l/2NA
3.8
where l is the wavelength of light and NA is the numerical aperture of objective lens. This relationship states that the shorter the wavelength, the greater is the resolution. Appropriate combination of magnification and resolution is defined by the NA of the lens system. When light passes through a minute object, it is diffracted into several wavefronts emerging at different angles from the incident light. This forms a cone of light originating from each point of the object, which is further gathered by the objective lens (Figure 3.8). The NA is, therefore, defined as the product of refractive index of the medium surrounding the object (n) and the sine of half of the angle of the cone of light (µ) that is accepted by the objective lens. NA = n (sin µ)
3.9
where n is the refractive index of the medium in which the object is placed and µ is half of the angle of the cone of light. Appropriate resolution is achieved by illuminating the specimen with adequate intensity of light that is uniform throughout the field of view.
Figure 3.8
Diagram depicting a cone of light originating from a point object
Microscopy
47
Often the visibility of image faces problem for uneven or inappropriate illumination which causes loss of resolution due to excessive glare or shadowing since the image of the light source falls on the same plane as the image of the specimen. In order to overcome such limitations, perfect illumination system was developed by August Kohler in 1893. It uses a condenser lens system for defocussing the image of the light source and illuminating the sample. A magnified image of the light source below the condenser, that is, at the aperture diaphragm, produces a wide cone of illumination required for optimum resolution of the specimen. The size of the condenser aperture diaphragm can be used to control the NA of the light cone that illuminates the sample and reduces unwanted light, producing a clear image with a good resolution. (iii) Contrast Transparent specimens lacking contrast are poorly imaged under a microscope. Contrast distinguishes the object from its background as well as from other adjacent features on the basis of the difference in their light intensity. Contrast is produced in the image due to various light phenomena such as absorption, scattering, diffraction, fluorescence, and colour variations. Staining of biological specimen is a usual method for enhancing the contrast. Contrast in an image is measured in terms of light intensity. Contrast per cent can be given by the following relation: Per cent contrast = [(Ib – Is)/Ib ] × 100
3.10
where Ib is the light intensity of the background and Is is the light intensity of the specimen. If the intensity of a specimen is less than that of the background, it appears darker and the contrast is referred to as positive contrast, while specimens that are lighter than the background are said to be in negative contrast.
3.5
PHASE CONTRAST MICROSCOPY The technique of choice for observing fixed and stained specimens always remained “transmitted brightfield microscopy,” since the objects visualized under this technique show significant absorption of various wavelengths of visible light on illumination, causing reduction in amplitude or intensity of light resulting into sufficient contrast. Obtaining contrast in unstained specimen, however, remained a question unanswered until the 1930s, when Frits Zernike discovered an optical technique which could translate the minute phase differences, unperceivable by unaided eyes, into corresponding amplitude
48 Bioanalytical Techniques changes, resulting into great contrast and making such objects easily visible. This technique revolutionized the microscopic visualization of living cells in detail and allowed the study of the dynamics of biological processes. Zernike was conferred upon the Nobel Prize in Physics for this discovery in 1953.
3.5.1
Development of Phase Contrast
When a phase object is illuminated through a beam of light, some of the light is diffracted by the object, while the rest of it does not deviate. The undeviated light rays (zeroth-order) pass through the specimen without interacting with any object and are referred to as “background waves,” and conversely the diffracted waves are produced as a result of interactions between the light and the specimen. Unstained objects, also called phase objects, appear transparent as they do not absorb light; however, they slightly alter the phase of the diffracted light, usually retarding it by approximately 1/4 of its wavelength, as compared to the direct or zeroth-order light passing through the specimen (Figure 3.9). This retardation is termed as “phase shift” and is caused by two factors: (i) the difference in refractive index between specimen and surrounding medium, and (ii) the thickness of the specimen. While observing under a microscope, both the background wave and the diffracted wave are focused by the objective at an intermediate image plane where they combine by the process of interference to produce a resultant wave. Since the amplitudes of background wave and resultant wave do not differ much, the image of phase objects does not acquire sufficient contrast so as to be visualized clearly, as is the case with unstained transparent biological specimen seen under brightfield microscopy (Figure 3.9a). In phase contrast microscopy, the phenomenon of phase shift is exaggerated by inserting a phase plate in the light path to speed up the zeroth-order background wave by 1/4th of wavelength. The direct undeviated light and the diffracted light waves, thus, become out of phase by 1/2 wavelength relative to each other (Figure 3.10). This results into destructive interference at the
Figure 3.9
Specimen under (a) brightfield microscope and (b) phase contrast microscope
Microscopy
Figure 3.10
Diagrammatic representation of phase shift in a light wave
Figure 3.11
Schematic representation of a phase contrast microscope
49
image plane, causing the image to appear darker against a lighter background with remarkable contrast (Figure 3.9b).
3.5.2
Phase Contrast Microscope
The arrangement of optical components in a phase contrast microscope is illustrated in Figure 3.11. The instrument consists of a light source, a condenser
50 Bioanalytical Techniques equipped with condenser annulus, a stage, objective lens, a phase plate, and an eyepiece. Light source is usually a tungsten halogen lamp or a mercury lamp which emits radiation in visible region. The incident light passes through a condenser annulus which is placed in the front focal plane of the condenser and is optically conjugate to the phase plate placed in the rear focal plane of the objective. This condenser annulus is made up of an opaque plate with a transparent annular ring, which is positioned in the front focal plane of the condenser that allows parallel light waves to illuminate the specimen. These light waves either pass through the specimen undeviated or are diffracted with a relative phase shift of 90 degrees. Both deviated and undeviated light waves are gathered by the phase objective equipped with a phase plate mounted at the objective’s rear focal plane. Depending on the objective lens magnification, that is 10X, 20X, or 40X, the size of the phase rings varies. Phase objectives are also matched to the appropriately sized annular diaphragm in the condenser by rotating the condenser to the suitable position. Objective lens of higher magnification require lesser width and diameter of phase plate but greater diameter of a condenser annulus. The phase plate segregates deviated and undeviated light waves emerging from the illuminated specimen at the rear focal plane of objective while altering the phase and amplitude of the surrounded (or undeviated) light wave. The etched part of the phase plate is narrow and optically thinner, and so when the undeviated light passes through the phase ring, it travels a shorter distance as compared to the diffracted light. This helps to speed up the direct undeviated zeroth-order light deliberately by 1/4 l. The net phase difference is between the direct undeviated light and the diffracted light now becomes 1/2 l, which is of paramount importance in creating contrast. The phase difference with which the direct light “speeds up” is calculated on a 1/4 l of green light; therefore, a green filter is placed in the light path which gives the best phase image with achromatic objectives as they are spherically corrected for green light. The etched ring part of the phase plate is also coated with an absorptive metallic layer to reduce the amplitude of direct light by 70–90 per cent and enhance the contrast.
3.5.3
Image Interpretation in Phase Contrast Microscopy
In positive phase contrast microscopy, a more common technique for biological viewing, the images appear darker on a bright background. The contrast in image is produced as a result of relative differences in intensity and amplitude of background or undeviated wave and the resultant wave emerging after the destructive interference at intermediate image plane (Figure 3.12).
Microscopy
Figure 3.12
51
Diagrammatic representation of undeviated wave, deviated wave, and resultant wave
The relative phase shifts enforced by various cellular/subcellular structures in the specimen depend upon their respective refractive indices and thickness. The structures having higher density will appear darker, such as nuclei, ribosomes, and mitochondria, while those having lesser density will appear lighter, such as vacuoles, vesicles, and cytoplasm. However, in phase contrast microscopy, there is no linear relationship between the intensity of an image and the density of the object. The contrast depends upon the absorption of light at the phase plate and the magnitude of phase shift. Phase images are often flawed with numerous optical artifacts such as shade-off and phase halos. Shade-off is seen in larger objects in which steadily reducing contrast from the centre of the object towards its edges makes it blur. Appearance of halo effect is a common feature in phase contrast microscopy seen as bright areas around the periphery of phase objects. This effect is prominent in specimen that induces large phase shifts, resulting into obscure details along boundaries. Despite these limitations, phase contrast microscopy is an ideal technique for studying thin specimens with high contrast and high resolution, enabling the examination of live cells in their natural state.
3.6 FLUORESCENCE MICROSCOPY Fluorescence microscopes were developed several decades after the discovery of the phenomenon of “Fluorescence” by Sir George G. Stokes in 1852. This technique revolutionized the research and studies in almost every field of science, such as cell biology, microbiology, plant and animal physiology, and organic and inorganic chemistry.
52 Bioanalytical Techniques Any molecule or cell that could fluoresce naturally or by binding to a fluorescing agent could become an object to fluorescence microscopy and revealed a whole new world of observations. Fluorescence tagging of living cells, labelling antibodies with fluorescent dyes, and probing of nucleic acids and proteins opened new gates to metabolic discoveries, enabling the study of intracellular dynamics such as rates of enzymatic reactions, signal transduction, DNA replication, etc. The technique has now become indispensable in analytical research, particularly in biological sciences, owing to its high sensitivity of emission profiles and high spatial resolution.
3.6.1 The Phenomenon of Fluorescence Sir George G. Stokes in 1852 first observed chemiluminescence in mineral fluorspar when it was irradiated with ultraviolet (UV) light, and he coined the term “fluorescence”. Such luminescence is exhibited by certain atoms or molecules which absorb light at a particular wavelength and subsequently emit light of longer wavelength; this wavelength shift was termed as “Stokes’ shift”. The phenomenon of photoluminescence can be explained with the help of Jablonski energy level diagram which depicts the electronic configuration of an atom (Figure 3.13). When an atom is irradiated with a short wavelength radiation, it absorbs a photon of light and upon excitation an electron of the atom is raised to a higher energy state (excited state). The excited electron immediately loses some energy in the form of heat and light emission, drops back to a lower excited vibrational state or singlet state, and finally returns to the ground state with emission of a longer wavelength of light. This emission is called fluorescence, which is very short lived having lifetime from 1 × 10 –5 to 10 –8 s. Sometimes
Figure 3.13
Energy-level diagram depicting electronic transitions giving rise to fluorescence and phosphorescence
Microscopy
53
the excited electrons make a forbidden transition to another lower excited state called triplet state and then to the ground state. In this process, the emission of radiation is considerably delayed, known as phosphorescence, which lasts longer having a lifetime of 1 × 10 –4 s to several minutes or even hours. For fluorescence imaging, the object should be either fluorescing in its natural form (termed primary or autofluorescence) or labelled with an extrinsic fluorescent molecule (known as secondary fluorescence). Autofluorescence is exhibited by numerous materials such as butter, chlorophyll, vitamins, minerals, crystals, resins, crude drugs, or even some organisms. But this fluorescence might not be specific or sufficiently intense so as to be utilized for analytical purpose. In this respect, the secondary fluorescence is far more important, which imparts specific fluorescence to cells, cell organelles, tissues, microorganisms, or macromolecules by labelling them with fluorochromes in a deliberate manner.
3.6.2
Fluorescence Microscope
The fluorescence microscope consists of a light source of high intensity which emits a wide spectrum of excitation wavelengths in UV–visible region. Mercury arc lamps are most commonly used in fluorescence microscope. Xenon arc lamps, light emitting diodes (LED), or laser systems can also be used selectively. The illumination passes through an excitation filter that selects appropriate radiation usually of a short wavelength and allows it to pass through while eliminating most of the other wavelengths (Figure 3.14). The selection of light source for fluorescence microscopy depends upon the absorption spectrum of fluorochromes and their quantum yield. Most fluorochromes are excited by shorter wavelength radiations such as UV, blue, and green. The excitation light further encounters the dichroic mirror through which it is reflected onto the objective to form an illumination cone, illuminating the specimen. The dichroic mirror is positioned in the light path at an angle of 45 degrees and is collinear with the objective and the specimen. It is designed to selectively reflect excitation wavelengths, while simultaneously transmitting short- and long-emitted wavelengths with high efficiency. In this way, the objective plays dual role; first it acts as a condenser and focuses the illumination of selected wavelength on the specimen, and second it gathers the fluorescence emitted from the excited specimen. The emitted radiation passes through an emission or barrier filter which transmits higher wavelengths, corresponding to yellow, orange, and red light. The excitation light of shorter wavelength reflecting from the specimen is able to pass through the dichroic mirror but is prevented effectively by the emission filter from reaching the detector. After passing through the emission filter, the emitted radiation is finally transmitted to the eyepiece or camera detection system.
54 Bioanalytical Techniques
Figure 3.14
3.6.3
Schematic representation of a fluorescence microscope
Fluorescence Microscope Images
As stated above, when the electrons drop from the excited state to the ground state, the loss of vibrational energy results into shifting of the emission spectrum to longer wavelengths than the excitation spectrum. This phenomenon is known as “Stokes’ shift” which allows the separation of excitation light from emission light. Use of appropriate filter ensures segregation of very weak light emitted through secondary emission by fluorochromes from intense excitation light. Figure 3.15 shows chromosomal translocation 9:22 in chronic myloid leukaemia in human using the technique of fluorescence in-situ hybridization.
3.7
ELECTRON MICROSCOPY As we have already discussed, the resolving power of a light microscope is limited by wavelength of light to 0.2 µm due to diffraction. The discovery of wave-like character of electrons in the 20th century led to the utilization of de Broglie wavelength for imaging as it was many orders of magnitude smaller (approximately 105 times smaller) than that of light, thus allowing imaging at atomic scales. A beam of electrons can be made to behave like a beam of electromagnetic radiation which is able to exhibit its wave properties. An electron accelerated under an electric field gains kinetic energy (E).
Microscopy
Figure 3.15
55
Fluorescence microscopic image showing chromosomal translocation 9:22 in chronic myloid leukaemia in human
E = eV = m0v2 /2
3.11
From this, velocity v of the electron can be derived as ____
2eV v = ____ m0
÷
3.12
where V is the acceleration voltage, e is the electron charge = –1.602 × 10 –19 C, and m0 is rest mass of the electron = 9.11 × 10 –31 kg. Momentum p of the electron is given by _______
p = m0v = ÷2m0eV
3.13
Now, wavelength l can be calculated from the de Broglie equation according to l = h/p _______
l = h/÷2m0eV where h is a Planck’s constant = 6.62 × 10 –34 J s.
3.14
56 Bioanalytical Techniques For example, if electrons are accelerated through a potential difference of 50 V, the wavelength obtained is about 0.17 nm. This wavelength is comparable to the atomic dimensions and can be utilized for imaging microstructures of organic and inorganic compounds. Electrons accelerated to a particular potential possess the same wavelength that constitutes a more or less monochromatic beam. The beam of high energy electrons interacts with the specimen within a certain area of cross section depending upon the electron beam energy. An accelerated electron penetrates into the electron cloud of an atom and its path is deflected, causing it to scatter at an angle. In some cases, even complete backscattering can occur, generating the backscattered electrons (BSE). Electron microscope utilizes all such interactions occurring between the matter and the highly accelerated electron incident on it. These electrons are focused with the help of electromagnetic lenses and their very short wavelength allows the specimen to be imaged with a very high spatial resolution as compared to the light microscope. The very first attempt to create an electromagnetic lens was made by Ernst Ruska and Max Knoll in 1931, after which they invented a transmission electron microscope (TEM). Eventually for this invention, Ruska was awarded the Nobel Prize in Physics in 1986, shared with H. Rohrer and G. Binnig for simultaneous development of scanning tunnelling microscope. Electron microscopy is now an indispensable imaging technique allowing resolution of the order of picometers and magnification of about 107 X that enables the study of topology, morphology, structural and compositional analysis of inorganic and organic materials, nano-materials, and lot more.
3.7.1
Matter–electron Interactions
Highly accelerated electrons are able to transmit through a thin specimen and are used for imaging in TEM. These electrons can be transmitted unscattered through the specimen, else they can be scattered elastically or inelastically. In elastic scattering, the incident electrons are deflected from their original path and scattered by atoms present in the specimen in an elastic fashion without any loss of energy. Such electron leaves the sample with almost same kinetic energy and velocity as it had initially. However, the trajectory of the electron may change after interaction with the specimen. Ee1 = E 0
3.15
Inelastic interactions involve transfer of energy from the incident electrons to the sample so that the energy of the electron gets reduced after interaction with the sample. Ee1 < Eo
3.16
Microscopy
57
These scattered electrons are then transmitted through the remaining portions of the specimen and can be collated using magnetic lenses to form a pattern of spots, each spot corresponding to a specific atomic spacing (a plane). This pattern can then yield information about the orientation, atomic arrangements, and phases present in the area being examined. Thick specimen does not allow incident electrons to transmit through, but instead results in interactions in same plane generating electrons with diverse energies. These electrons are used for imaging in scanning electron microscopy. Some of the electrons and radiations generated as a result of matter–electron interaction and used in electron microscopy imaging are discussed below. (i) Backscattered electrons Electrons reflected back after colliding with the specimen by elastic scattering are called BSE. These electrons can be distinguished on the basis of their high kinetic energy as they escape from the sample with energies quite equivalent to the primary-beam energy. The fraction of primary electrons that escape as BSE is given as the backscattering coefficient (hb) which varies directly with the specimen’s atomic number. nBSE Number of electrons backscattered ___________________________________________ 3.17 hb = ____ nB = Number of beam electrons incident on specimen where nBSE is the number of BSEs and nB is the number of beam electrons incident on the specimen. The higher atomic number elements cause more production of BSEs and appear brighter than lower atomic number elements. Therefore, BSE images show contrast in different parts of the specimen having different average atomic numbers. (ii) Secondary electrons During inelastic interactions, the incident electrons may transfer their energy to the outer shell or valence electrons. These electrons are electrostatically weakly bound to the nucleus and need only a small amount of energy to overcome the binding force and to eject out into the vacuum. These are called secondary electrons (SEs). Most SEs eject with a very small kinetic energy ( K+ > NH4+ The order of displacement effectiveness for commonly used anions: PO43– > SO42– > COO – > Cl– These rankings correlate with the Hofmiester series but it has to be borne in mind that the strongest eluting salt is not always best. Ideally, several salts should be tested, and finding optimum elution conditions often involves trial and error.
6.4.4.1
Step and gradient elution
Figure 6.24 depicts ion-exchange separation wherein proteins are eluted by increasing the ionic strength of a buffer (typically with NaCl) using linear gradient or step elution. The UV absorbance and conductivity traces show the elution of the protein peaks and the changes in salt concentration, respectively, during elution (Figure 6.24a). In gradient elution, a linear increase in salt concentration of elution buffer is brought about using special gradient mixers. In step elution, the increase in salt concentration is brought about in discreet steps (Figure 6.25b).
6.4.4.2
Regeneration
This step simply involves removing all bound protein from the stationary phase so that it is ready for another process. Key considerations for choosing ion-exchange system are as follows: (i) Stability of proteins or analytes (a) When the analyte is stable below the pI value, choose a cation exchanger. (b) When the analyte is stable above the pI value, choose an anion exchanger. (ii) If the molecular weight is lower than 10,000, then we choose the smaller size of the matrix particles. When the molecular weight is greater than 10,000, then the bigger size of the matrix particles can be used.
6.4.4.3 Applications of ion-exchange chromatography Ion-exchange chromatography has the following applications: (i) Preparation of high purity water, water softening, and decontamination; for example, water softening is accomplished by exchanging calcium Ca2+ and magnesium Mg2+ cations against Na+ or H+ cation. Another application
Chromatographic Techniques
Figure 6.24
(ii) (iii)
(iv)
(v)
169
Elution in ion-exchange chromatography: (a) linear gradient (b) step gradient
for ion exchange in domestic water treatment is the removal of nitrate and natural organic matter. Purification of charged molecules such as proteins, amino acids, and DNA or RNA. Separation of uranium from plutonium and other actinides, and then plutonium and uranium are available for making nuclear energy materials, such as new reactor fuel and nuclear weapons. Separation of sets of very similar chemical elements, such as zirconium and hafnium ion exchangers are used in nuclear reprocessing and for the treatment of radioactive waste. Ion-exchange resins in the form of thin membranes are used in chloralkali process, fuel cells, and vanadium redox batteries.
170 Bioanalytical Techniques 6.4.5
Gel Permeation Chromatography
Gel permeation chromatography (GPC) is also known as gel-filtration chromatography, molecular sieve chromatography, and size exclusion chromatography (SEC). Gel-filtration chromatography separates molecules according to their size and shape. A column of uniformly sized microparticles of cross-linked copolymers, usually made of either styrene or divinylbenzene, makes the stationary phase. The pore size of the gel is controlled within a narrow range. Therefore, the analytes that are larger than the pores will be completely excluded from the pores and would travel only through the interstitial spaces and appear first in the eluate. The smaller analytes will be distributed within the mobile phase inside and outside the particles and will, therefore, emerge later.
6.4.5.1
Gel-filtration matrices
A wide variety of gel-filtration matrices are available commercially, such as Sephadex (dextran beads), Sepharose, both from Pharmacia and Bio-Gel A (agarose), and Bio-Gel P (polyacrylamide) from Bio-Rad Labs. Alternatively, other materials such as polyacryloylmorpholine and various polystyrenes have also been used. These materials can be made into microparticulate beads with different degrees of porosity and, hence, will fractionate different size ranges of proteins. Dextran-based gels allow fractionation of macromolecules up to about 800,000 molecular weight whereas agarose gels, because of their greater porosity, can be used for macromolecules up to several million molecular weight. Recently cross-linked dextran gels (Sephacryl) have been marketed which are exceptionally rigid and stable matrices for fractionation of proteins. New types of commercial gels are constantly being developed. Table 6.10 summarizes the properties of the commonly used matrices.
6.4.5.2
Key steps in gel-filtration chromatography
Gel-filtration chromatography has the following steps: (i) Swelling of gel The gel-filtration beads need to be swollen in appropriate buffer (sometimes these are available as pre-swollen matrix from the manufacturer). Most commonly used buffers are Tris–HCl, sodium phosphate, and sodium acetate. An ionic strength of at least 0.05 M is used to reduce non-specific interactions between the proteins being separated and the chromatographic matrix. (ii) Column packing The gel beads, present as slurry in the chosen buffer, are poured into a glass or plastic chromatography column of suitable dimensions and allowed to settle by gravity taking care not to introduce
Chromatographic Techniques
171
Table 6.10 Characteristics of gel-filtration media used in this simulation Matrix name (trade Bead type names)
Approximate fractionation range for peptides and globular proteins (molecular weight)
Sephadex G-50
Dextran
1,500–30,000
Sephadex G-100
Dextran
4,000–150,000
Sephacryl S-200 HR
Dextran
5,000–250,000
Ultrogel AcA 54
Polyacrylamide/agarose
6,000–70,000
Ultrogel AcA 44
Polyacrylamide/agarose
12,000–130,000
Ultrogel AcA 34
Polyacrylamide/agarose
20,000–400,000
Bio-Gel P-60
Polyacrylamide
3,000–60,000
Bio Gel P-150
Polyacrylamide
15,000–150,000
Bio-Gel P-300
Polyacrylamide
60,000–400,000
air bubbles. Peristaltic pumps can be used to pack the gel chromatography columns. After washing and equilibration of the column with buffer alone, the protein mixture in the buffer is applied to the top of the column and the eluate is collected at the column base in a series of fractions. A layer of buffer is always kept in contact over the column bed to prevent the gel from drying. (iii) Washing the column The column is washed with 2 to 3 bed volumes of buffer using a peristaltic pump. This ensures appropriate equilibration of the column with the buffer. The flow rates in gel permeation chromatography should never exceed the limits given by the bead manufacturer since the beads may deform at higher pressure. (iv) Application of sample The sample dissolved in appropriate buffer is applied without disturbing the gel. (v) Elution Elution is performed by using appropriate buffer and fractions are collected manually or by using automatic fraction collector. The fractions collected are analysed for the desired compounds. For example, the proteins can be monitored using UV absorption at 280 nm.
6.4.5.3
Determination of molecular weight of macromolecules using gel filtration chromatography
Gel-filtration chromatography can be used for determining the molecular weight of macromolecules. A column can be calibrated with standard proteins of known molecular weight. A sample macromolecule mixture is applied to the same column; measurement of the elution volume of the protein of interest
172 Bioanalytical Techniques allows its molecular weight to be deduced by reference to the standard curve. However, it should be noted that the shape of the protein molecules also plays an important role in gel filtration. Long, extended polypeptides and proteins tend to behave as though they were larger, globular protein molecules. Therefore, a calibration curve is only as accurate as the nature of the protein standards used for constructing it will allow. The void volume Vo of the gel permeation column is the unoccupied volume of a gel permeation column. It is determined by using a very large-sized molecule (such as Blue dextran molecular weight 2 × 106 Daltons) that cannot penetrate the pores of the gel matrix but remains totally unretained. Therefore, the void volume is determined as the volume of liquid collected from the column from the moment the sample is applied to the column and its elution from the column. To determine the molecular weight of an unknown molecule, a large molecular weight compound such as dyed dextran is applied to the size exclusion chromatography column along with the standards (molecules of known molecular weight) and the sample. The volume at which blue dextran elutes is the “void volume”. Similarly, elution volume of standard molecules of known molecular weight is recorded and a calibration curve is drawn. A plot of molecular weight of the standard proteins versus Vr /Vo gives a calibration curve for the column (Figure 6.25). As the proteins pass down the column, they penetrate the pores of the gel beads to different extents and so travel down the column at different rates. All proteins which exceed the maximum size of the pores will not be able to enter the beads. These proteins will, therefore, be distributed only in the solution between the beads and the elute from the column first in the “exclusion volume” (also called void volume). Kd =
Concentration in mobile phase ______________________________ Concentration in stationary phase
6.26
To compare data between columns, a relative elution volume is used. In SEC, this ratio can be expressed as a ratio of retention times or retention volumes. Kd =
V r – Vo _______ Vs
where Vr is the retention volume of the protein, Vo is the void volume, and Vs is the volume of the stationary phase.
6.27
Chromatographic Techniques
Figure 6.25
173
Calibration curve for determining of molecular weight of proteins using gel permeation chromatography
Note Proteins of known molecular weights: (a) Ferritin (b) Catalase (c) Ovalbumin, and (d) Myoglobin were applied as standards to gel filtration column along with a protein of unknown molecular weight. The volume of retention of the standard and unknown proteins were recorded and plotted in a semilog plot against Vr/Vo. A straight line graph is obtained and the molecular weight of the protein can be found out using the equation of the straight line thus obtained.
However, Vs or the stationary phase volume is difficult to calculate accurately. Instead, the average distribution constant used is Kavg = [Vr – Vo ]/[Vt – Vo ]
6.28
where Vt is the maximum retention volume experienced by a small molecule such as a dye labelled amino acid. This is a close approximation of the stationary phase volume in the column. Kavg values have been determined for a number of gels and for a number of proteins.
6.4.5.4 Applications and advantages The primary application of gel-filtration chromatography is the fractionation of proteins and other water soluble polymers. SEC is a widely used technique for the purification and analysis of synthetic and biological polymers, such as proteins, polysaccharides nucleic acids, and other molecules. The main advantage of gel-filtration chromatography is that the analytes do not bind to the gel-filtration matrix. This permits the separation at operating conditions suitable to the analytes than to the conditions required to promote binding/elution operations as in ion-exchange or affinity chromatography (AC).
174 Bioanalytical Techniques The separations can be performed with buffers containing stabilizing agents and at low temperatures. The eluents from ion-exchange chromatography can be directly applied to gel filtration without the need of desalting. It has to be noted, however, that gel chromatography leads to dilution of sample. Therefore, it is beneficial to apply a concentrated sample.
6.4.6
Hydrophobic Interaction Chromatography
Hydrophobic interaction chromatography (HIC) is a technique used for separating peptides, proteins, and other biological molecules based on their degree of hydrophobicity. The HIC mobile phase comprises a high concentration of salting-out agent, typically ammonium sulfate, which increases the hydrophobic interaction between the solute and the stationary phase. Most proteins, and to a much lesser extent hydrophilic molecules, for example, DNA and carbohydrates, have hydrophobic areas or patches on their surface. Solvation of these patches is energetically unfavourable and results in the formation of hydrophobic cavities in the aqueous mobile phase. The promotion of the hydrophobic effect (by addition of lyotropic salts) drives the adsorption of hydrophobic areas on a protein to the hydrophobic areas on the solid support. This is thermodynamically favourable in that it reduces the number and volume of individual hydrophobic cavities. Reduction of hydrophobic interaction by decreasing the concentrations of lyotropic salts results in desorption from the solid support. HIC is unique in that the proteins bind at high-salt concentration and elute at low-salt concentration. Hydrophobic interaction chromatography and reversed-phase chromatography are closely related techniques. Both are based upon interactions between solvent-accessible non-polar groups (hydrophobic patches) on the surface of the solute and the hydrophobic ligands of the stationary phase. In practice, however, HIC and reversed-phase chromatography differ since reversed-phase stationary phases are more highly substituted with hydrophobic ligands than HIC stationary phases. Since protein binding in reversed-phase chromatography (RPC) is usually very strong, it requires polar solvents for elution whereas HIC allows working with non-polar, non-denaturing mobile phases. This is an advantage as such solvents can denature the proteins.
6.4.6.1
HIC matrices
Matrices for HIC are commonly agarose based that are substituted with hydrophobic groups. Generally, linear alkanes (butyl, octyl) that show pure hydrophobic nature are used for substitution. Alternatively, aromatic (phenyl) ligands that provide some degree of aromatic interactions are also used. An
Chromatographic Techniques
175
increase in alkyl chain length and ligand density increases the binding capacity. However, higher degree of binding capacity makes the elution of proteins difficult.
6.4.6.2
Key steps in HIC
The chromatographic medium is first activated to expose the hydrophobic groups. Binding of sample is performed at high-salt concentrations, which helps in exposing the hydrophobic groups on the protein for interaction with the matrix. The column is washed to remove the unbound material. The elution is performed with successively decreasing salt concentration (linear or step gradient may be used).
6.4.6.3 Advantages of HIC The following are the advantages of HIC: (i) A large volume of samples can be loaded; therefore, purification steps that generate large sample volume can be coupled with this method. (ii) Samples with high ionic strength can be used; therefore, HIC is good for samples after ammonium sulfate fractionation. (iii) The sample is eluted with low salt. HIC is well suited for being used before gel filtration, ion exchange, and AC. These techniques usually require an additional step of ionic strength. However, when these are performed after HIC, this step may be eliminated. (iv) A sample can be used in the ion-exchange chromatography step.
6.4.7 Affinity Chromatography Affinity chromatography separates molecules based on the reversible interaction between target protein and the specific ligand attached to a chromatography matrix. The interaction can be specific (such as antibodies binding to protein A or protein G) or non-specific (for example, a histidine-tagged protein binding to metal ions). Many other types of biological interactions can be utilized in AC, including enzyme–substrate, protein–protein, or nucleic acid–protein interactions. Affinity chromatography is unique in purification technology since it is the only technique that enables the purification of a biomolecule on the basis of its biological function or individual chemical structure. The technique offers high selectivity, hence high resolution, and usually high capacity for the protein(s) of interest. Purification can be in the order of several thousand-fold and recoveries of active material are generally very high. Purification that would otherwise be time consuming, difficult or even impossible using other techniques can often
176 Bioanalytical Techniques be easily achieved with AC. This technique can be used for separating active biomolecules from denatured or functionally different forms, isolating pure substances present at low concentration in large volumes of crude sample, and removing specific contaminants. Table 6.11 gives the different types of interactions exploited for separation with AC:
6.4.7.1 Affinity chromatography matrices Gel exclusion stationary phase matrices work well in AC because they possess the desirable properties such as: (i) physically and chemically stable under most experimental conditions, (ii) relatively free of non-specific adsorption effects, (iii) very large pore sizes, and (iv) reactive functional groups for ligand attachment. The examples of matrices used are agarose (for example, Sepharose 4B), polyvinyl, polyacrylamide, and controlled porosity glass.
Table 6.11 Affinity chromatography ligands Analyte
Ligand
Examples
Enzyme
Substrate, inhibitor, substrate analogue, inhibitor analogue
Serine proteases: Benzamidine (inhibitor analog)
Antigen
Antibody
Diagnosis of pathogens
Antibody
Protein A or G
Isolation of antibodies
Hormone receptor
Hormone
Thyroid hormone receptor - to this T3 affinity
Lectin
Polysaccharide
Resolution of mannans
Nucleic acids
Complementary base sequence, histones, nucleic acid polymerase, nucleic acid binding protein
Glutathione-Stransferase or GST fusion proteins.
Glutathione
Proteins
Metal ions (Ni and Cu)
Biotinylated substances
Streptavidin
Poly (His) fusion proteins, native proteins with histidine, cysteine and/or tryptophan residues on their surfaces.
Chromatographic Techniques
177
6.4.7.2 Attachment of ligand to the matrix The ligand is covalently attached to the AC matrices. The matrix is activated so that a functional group that can covalently react with the ligand is introduced and the ligand is attached. The covalent attachment of the ligand to the matrix requires prior activation of the matrix. Many activated matrices are also available commercially: for example, CNBr-activated agarose, Sepharose 4B, 6-aminohexanoic acid (CH)-agarose, and epoxy-agarose.
6.4.7.3
Key steps in affinity chromatography
The “column equilibration” and “sample application” are similar to the gelfiltration chromatography. After application of the sample, the column is washed with a “washing buffer”, so that the components of the sample mixture which do not bind specifically to the column are washed out of the column.
6.4.7.4
Elution in affinity chromatography
Several strategies can be taken up for elution of bound analytes in AC as described in the following section: (i) pH elution A change in pH alters the degree of ionization of charged groups on the ligand and/or the bound protein. This change may affect the binding sites directly, reducing their affinity, or causing indirect changes in affinity by alterations in conformation. Usually, the a decrease in pH is brought about in a step-wise manner for elution of bound solutes. The chemical stability of the matrix, ligand, and target protein determines the limit of pH that may be used. (ii) Competitive elution Selective eluents are often used for separating substances on a group specific medium or when the binding affinity of the ligand or target protein interaction is relatively high. The eluting agent competes either for binding to the target protein or for binding to the ligand. (iii) Ionic strength elution Increased ionic strength (usually NaCl) is applied as a linear gradient or in steps. (iv) Elution dependent on reduced polarity of solvent Conditions are used to lower the polarity of the eluent to promote elution without inactivating the eluted substances. Dioxane (up to 10 per cent) and ethylene glycol (up to 50 per cent) are typical of this type of eluent. (v) Chaotropic agents Chaotropic agents, such as guanidine hydrochloride or urea, can be used that alter the structure of the attached proteins. However, chaotropes should be avoided if isolation of the biological activity of the analyte is desired since it is likely to denature the eluted protein.
178 Bioanalytical Techniques 6.5
PROTEIN PURIFICATION STRATEGIES The goal of any separation technique is to obtain a high yield of highly pure and active protein in minimal number of steps. Since biological samples (fermentation broths, cell extracts, serum, etc.) contain many impurities that possess properties similar to the desired molecule, the final purity is not obtained in a single step. A train of purification steps has to be performed before the final purity is achieved (Figure 6.26). However, this warrants high recovery at each step because yield would reduce with every step. For example, a purification scheme having four steps with 80 per cent step yield will give 41 per cent final yield, while four steps at 60 per cent step yield will give a 13 per cent final yield. While designing a purification strategy, it is advisable to use the steps having maximum selectivity as early as possible in a purification sequence. Generally, higher selectivity is afforded by the techniques that utilize unique properties of the desired molecules for achieving separation; for example, AC. Moreover, it is understandable that the fewer the steps, the faster the preparation, the lower the protein losses, and the lower the cost of the purification procedure. Some of the key considerations in designing a purification procedure are: (i) to have a convenient assay to follow purification; (ii) to have a concentrated starting material; (iii) to have precautions for minimization of product inactivation; (iv) to minimize number of steps; (v) to avoid unnecessary dialysis and delay; and (vii) to have high resolution. Generally, the purification of proteins involves the following sequence: (i) Precipitation (ii) Ion exchange (iii) Affinity (iv) Gel filtration
Figure 6.26
A possible network of chromatographic schemes
Chromatographic Techniques
179
To exemplify a typical protein purification sequence, a hypothetical protein (enzyme) was purified sequentially using the steps given in Table 6.12. To determine the success of the protein purification scheme, the specific activity of the protein obtained at each step was determined. Also the product obtained at each step was subjected to SDS-PAGE analysis (Figure 6.27). At each step, the following parameters were measured: (i) Total protein The quantity of protein present in a fraction is obtained by determining the protein concentration of a part of each fraction and multiplying by the fraction’s total volume. (ii) Total activity The enzyme activity for the fraction is obtained by measuring the enzyme activity in the volume of fraction used in the assay and multiplying by the fraction’s total volume. (iii) Specific activity This parameter is obtained by dividing the total activity by the total protein. (iv) Yield This parameter is a measure of the activity retained after each purification step as a percentage of the activity in the crude extract. The amount of activity in the initial extract is taken to be 100 per cent. (v) Purification level This parameter is a measure of the increase in purity and is obtained by dividing the specific activity, calculated after each purification step, by the specific activity of the initial extract. As we see in Table 6.12, the first purification step, ammonium sulfate precipitation, led to an increase in purity of only 2.5 fold but recovery was 75 per cent. After ammonium sulfate precipitation, the sample was dialysed to remove the salt and the dialyzed sample was applied to an ion-exchange column. At this step the purification achieved was 9.6 fold but the yield was reduced to 70 per cent. The gel-filtration chromatography led to an increase of purity to 117 fold but the yield further fell to 47 per cent. The final step used was Table 6.12 Summary of a hypothetical purification scheme of an enzyme Step
Total Total activity Specific activity, Yield (per Purification protein (mg) (units) (units mg–1 cent) level
Cell disruption
15,120
150,120
10
100
Ammonium sulfate precipitation
4,500
11,2500
25
75
2.5
Ion-exchange chromatography
1,100
105,500
96
70
9.6
Gel-filtration chromatography
60
70,000
1,167
47
117
50,400
28,000
34
2,800
Affinity chromatography
1.8
1
180 Bioanalytical Techniques
Figure 6.27 Expected electrophoretic pattern for a hypothetical protein purification as discussed above. Lane 1: After cell disruption; Lane 2: Ammonium Sulphate precipitation; Lane 3: Ion Exchange chromatography; Lane 4: Gel filtration; Lane 5: Affinity Chromatography purified product.
AC that resulted in a purification level of 2800 fold. The yield obtained was 34 per cent. Therefore, more than 50 per cent of our product had been lost during purification. The SDS PAGE analysis was performed on the samples obtained after each step. It can be seen that the number of bands decreases with the level of purification (Figure 6.27). A good purification scheme takes into account both purification levels and yield. A high degree of purification and a poor yield leave little protein with which to experiment. A high yield with low purification leaves many contaminants (proteins other than the one of interest) in the fraction and, hence, complicates the interpretation of experiments.
7 Electrophoresis
7.1 INTRODUCTION It is truly mesmerizing to see biomolecules such as DNA and proteins running into the gel in the process of electrophoresis. Since many decades electrophoresis has been a method of choice for separation of charged molecules under the influence of electric field. In the 20th century, W. B. Hardy realized that many biologically important molecules, such as enzymes and other proteins, displayed characteristics of electrophoretic mobility as they existed in ionized state in solution. In subsequent years, this property was explored for establishing this indispensable technique based on migration of charged molecules through a conductive medium under the influence of an externally applied electric field. Several biochemists utilize the property of electrophoretic mobility in a predictable manner to determine molecular weight of nucleic acid fragments and isoelectric points (pI) of many enzymes. The history of electrophoresis dates back to the contribution of Arne Tiselius in the 1930s. He developed a procedure and apparatus for moving boundary electrophoresis of proteins in a buffer solution for which, along with his work on adsorption analysis, he received the Nobel Prize in Chemistry in 1948. The apparatus that Tiselius developed is still used for separating proteins in solution by electrophoresis and for analytical purposes such as determination of pI, molecular weight, and related physical properties. Majority of the early electrophoretic studies of biomolecules were carried out using moving boundary electrophoresis in liquid phase. The major disadvantage in this technique was that the separating boundaries of migrating species overlapped and did not separate into clear zones due to convection currents in the medium and lateral diffusion of ions. To overcome this limitation, in further
182 Bioanalytical Techniques years the electrophoretic technique was refined by the use of an inert support medium saturated with a buffer which holds the molecules as they migrate in the conducting medium, thus preventing convection and diffusion. In 1939, Dionysius Von Klobusitzky and P. Konig successfully applied electric field to paper strips saturated with an electrolytic solution to separate the components of snake venom. Free solution and paper electrophoresis are further discussed in Sections 7.3 and 7.4. The use of solid or semi-solid supporting media such as filter paper and gels gave way to zone electrophoresis. This technique offered two main advantages: first, the sample requirement decreased several times and second, it stabilized the migration of molecules reducing their diffusion resulting into markedly clear separation bands. Zone electrophoresis has now become a basic separation tool in molecular biology laboratories and is widely used for separating macromolecules carrying charge such as amino acids, peptides, proteins, and nucleic acids.
7.2
PRINCIPLES OF ELECTROPHORESIS In the process of electrophoresis, migration of particles or molecules takes place through a fluid medium (mobile phase) which is a buffer of definite pH or a pH gradient. Such migration is characterized by the term “electrophoretic mobility”, and is defined as the measure of the rate of migration of molecules under applied electric field. Depending upon the pH of the buffer used, either cationic or anionic species predominate, which on application of electric field flow in corresponding direction establishing an electro-osmotic flow (EOF). This bulk flow tends to sweep all molecules in that direction. The rate of migration of molecules depends not only upon their relative charge, mass, and conformation, but also on other factors such as pore size of gel, types of ions in the buffer system, and temperature. Thus, the net rate of migration achieved by the molecules becomes the very basis of their separation under electrophoresis.
7.2.1
Concept of Electrophoretic Mobility
The influence of electric field on a molecule depends upon its effective charge and can be explained by the “double layer theory”. The molecules, bearing some surface charges, are immersed in a fluid medium and are screened by a diffuse layer made of counter charges. The diffuse layer is formed under the influence of electric attraction and thermal motion of free ions in the surrounding fluid. The net electric charges in the surface layer as well as in the diffuse layer are opposite and equal in magnitude, resulting in a neutral structure.
Electrophoresis
Figure 7.1
183
Forces acting on a molecule during its migration under electric field
This interfacial double layer shields the molecule from the field and, therefore, the electric field interacts with the particle indirectly. Under the force exerted by the electric field, the double layer is pushed in the direction of electric field which further pulls the particle along with it. During this movement, the molecule experiences three kinds of forces: electrostatic force, electrophoretic retardation force, and frictional force (Figure 7.1). When a potential difference is applied across the electrodes, it generates a potential gradient (E), and the molecules suspended in the medium, carrying surface charge (q), experience an electrostatic or Coulomb force which drives a charged molecule towards an electrode. This force is given as Fel = qE
7.1
The force exerted by the electric field on the diffuse layer is opposite in direction as compared to that exerted on the surface charge of the molecule. The ions in diffuse layer transfer some of this force to the particle, called electrophoretic retardation force (Fret), which is lesser in magnitude and opposes the movement of the molecule. The movement of the molecule in fluid medium is also resisted by frictional force (Ff) generated by the Stokes’ drag, which is proportional to the velocity (v) of the molecule and is given as Ff = – fv
7.2
where f is the frictional coefficient. f is given as the following relation f = 6phrs
7.3
184 Bioanalytical Techniques where h is the viscosity of the medium and rs is the Stokes’ radius of the molecule. The magnitude of frictional resistance depends upon the hydrodynamic size and shape of the molecule, pore size of the electrophoretic medium, and the viscosity of the buffer. The larger value of rs leads to the larger frictional coefficient and the larger Ff which resists motion towards the electrode. The electrostatic force is balanced by the other two counteracting forces and is represented as Fel = – (Fret + Ff )
7.4
For a weak electrical double layer, the electrophoretic retardation force can be ignored; therefore Fel = (–Ff)
7.5
From Equations 7.1 and 7.2, we get n = qE /f
7.6
The rate of movement of particle through the fluid is proportionately dependent upon the applied electric field strength and can be expressed as electrophoretic mobility (µ). m = n/E
7.7
From Equations 7.3, 7.6, and 7.7, we get m = q/6phrs
7.8
Therefore, macromolecules of different charge density can be separated by electrophoresis. The electrophoretic mobility µ is defined as the rate of migration (in cm/s) per unit electric field strength (V/cm) of a charged particle in electrophoresis. It is proportional to the charge density (q/rs) of the particle.
7.2.2
Behaviour of Biomolecules in Electrophoretic Matrix
Biological macromolecules exhibit dynamic behaviour in electric fields involving electrostatic interactions with cations and anions existing in surrounding buffer and with those present on the surface of electrophoretic matrix. Most common biomolecules resolved electrophoretically are proteins and nucleic
Electrophoresis
185
acids. Proteins are polymers of amino acid units joined by peptide bonds and arranged in complex conformations of variable degree such as primary, secondary, tertiary, or quaternary. These are amphoteric molecules and their net charges are determined by free amino group, carboxyl group, and ionizable alkyl groups attached to it. Proteins bear positive charge in acidic solution and a negative charge in basic solution, and at intermediate pH value they carry zero net charge, which is called isoelectric point of a protein (Figure 7.2a). This pH dependence of charge affects their mobility in terms of magnitude and direction of migration. Nucleic acids such as DNA and RNA are formed by polymerization of individual nucleotide units, connected via phosphodiester linkages. The phosphate groups present on sugar–phosphate backbone are acidic in nature and tend to lose a proton in solution so that it attains net negative charge (Figure 7.2b). The polar groups present in proteins and nucleic acids render some net charge to these polymers (as shown in Figure 7.2a), which allows them to migrate in the applied electrical field in a predictable manner. The greater the net charge, the greater the mobility of the molecule towards the opposite pole. The net charge of a particular molecule depends upon the pH of the surrounding buffer. If all the molecules bear same amount of charge, the rate of migration is theoretically said to be directly proportional to their sizes. In electrophoresis, macromolecules resolve on the basis of this phenomenon. But the mobility of the biomolecules in an applied external electric field is governed by several factors other than size. Proteins, for example, are incredibly diverse and complex in terms of size, shape, and conformations. The thermal and electric field stress may introduce conformational alterations in proteins and nucleic acids, and bring about unexpected changes in their mobility. Thermal fluctuations generated by the electric field also tend to change the dipole moment of the molecules. Under this fact, relationship between the rate of migration of molecules under the influence of electric field and their size does not remain linear. While carrying out electrophoresis using semi-solid matrices, the mobility of macromolecules also depends upon concentration and viscosity of the gel. Smaller pores and higher viscosity slow down the migration rate. Due to its larger pore size, agarose gel is more suited for the electrophoretic separation of nucleic acids (200 kDa). In the case of DNA, polyacrylamide can be used for separating fragments smaller than about 500 bps. Typically polyacrylamide gels are used within the range of 3–30 per cent which gives approximate range of pore size from 200 nm to 20 nm. Procedures requiring sieving effect and separation of proteins on the basis of size are carried out in gels having concentration 12 per cent or higher. Proteins with molecular weights ranging 10–100 kDa may be separated in 8–12 per cent polyacrylamide gels. Low-concentration gels (3–4 per cent) are used for electrophoresis of proteins requiring free mobility through the gel as in IEF, stacking gel of discontinuous SDS PAGE, electrophoresis of larger molecules such as nucleic acids, or sequencing gels. (iv) Cellulose acetate Cellulose acetate electrophoresis is an accurate and a simple method for protein quantification. It is carried out in pre-cast thin plates or strips. Electrophoretic mobility of proteins in cellulose acetate electrophoresis is determined by their charge. Migration takes place on Table 7.1 Acrylamide concentrations suitable for separation of proteins in different molecular size range Acrylamide in resolving gel (in per cent) Uniform gel
Gradient gel range
Separation size range 5
35–200 kDa
7.5
25–200 kDa
10
15–200 kDa
12.5
15–100 kDa
15
10–100 kDa
5–15
10–200 kDa
5–20
10–200 kDa
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197
the buffer film on the surface of the cellulose acetate layer. These strips require very less amount of sample to be loaded (0.25–0.50 µl), lesser stain volume, and shorter gel running time (15–25 min). Cellulose acetate medium can be used in dry or wet (humid) forms. Dry cellulose acetate strips offer isotropic, highly stable, and integral matrix for ideal flow conditions. These have high porosity and low adsorption capacity exhibiting minimum diffusion; therefore, they can also be used for high voltage electrophoresis. Dry strips can separate fairly wide range of materials that migrate in an electrical field but the concentration of the sample mixture to be separated should be more than 1 mg/ml. Wet strips are thoroughly impregnated with buffer for approximately 1 h before use. These have unique conical pores which reliably absorb the samples and allow optimal application for large volume samples without streaking. These strips can give sharp bands under low voltage electrophoresis also. They are extremely resistant to strong alkalis, therefore wide range of dyes including alkaline dyes can be used for staining the sample. Most common immunological application of cellulose acetate electrophoresis is analysis of haemoglobin for diagnosis of haemoglobinopathies such as sickle cell anaemia and thalassemia. In this technique, the haemoglobin is treated under alkaline conditions so that it attains negative charge. However, different variants of haemoglobin carry different net charge, therefore, migrate in the gel at different speeds. On the basis of their respective rate of migration, they can be compared with known standards and can thus be classified. Cellulose acetate is an excellent medium for isoenzyme analysis such as lactate dehydrogenase (Figure 7.10), malate dehydrogenase, phosphokinase, and so on for identification of genotypes and factors affecting genetic variation in populations. Besides, it is also used for the separation of serum proteins for monoclonal gammapathies, lipoproteins, glycoproteins, mucopolysaccharides, cerebrospinal fluid, urine, and other body fluids.
7.5.2 7.5.2.1
Gel Conditions Denaturing or non-denaturing system
The mobility of molecules in gel electrophoresis is influenced by both charge and size, so different protein molecules with similar molecular weights may migrate differently owing to their characteristic charges. Moreover, the mobility is also affected by the hydrodynamic size of proteins and their secondary, tertiary, or quaternary conformations. Due to these reasons, the measurement of molecular weight or conformational analysis of the biological molecules
198 Bioanalytical Techniques
Figure 7.10
Cellulose acetate electrophoresis patterns of LDH isozymes in human serum for detection of liver disease
is not always feasible by general electrophoresis technique. Several protein electrophoretic studies require prior denaturation of the molecules so that all the molecules of the protein attain an unfolded rod-shaped conformation to become identical in shape and have a uniform charge. The electrophoretic mobility of these molecules will now depend essentially on their size. The most common agent used for dissociating proteins is an anionic detergent SDS, along with a thiol reagent such as b-mercaptoethanol (b-ME) or dithiothreitol (DTT). SDS binds proportionately to protein molecules (1.4 g of SDS binds per gram of protein), denatures them to their primary structure by disrupting non-covalent bonds in the proteins, and coats them with uniform negative charge, irrespective of their intrinsic charges (Figure 7.11). The protein sample is often heated about 100°C which further promotes protein denaturation and helps SDS in binding to the protein. Reducing agent (b-ME
Figure 7.11
Conversion of protein molecule to linear polypeptide of uniform negative charge by SDS
Electrophoresis
199
or DTT) assists in the protein denaturation by reducing all disulfide bonds and breaking tertiary and quaternary structures. The SDS protein complexes now have essentially identical charge densities and migrate in polyacrylamide gels strictly according to the polypeptide size. Urea gradient gels are also used for electrophoresis involving denaturantinduced dissociation of proteins and nucleic acids. These are commonly used where the electrophoretic mobility of proteins and nucleic acids need to depend exclusively on their sizes. Urea brings about dissociations by disrupting hydrogen bonds making DNA single-stranded and rendering more flexibility in its movement. For complete denaturation of proteins, it is usually combined with thiol reagent b-ME. Urea is non-ionic so it can be incorporated in gels without affecting the intrinsic charge of biomolecules so that their separation occurs on the basis of both size and charge. However, the combination of size and charge fractionation may prevent accurate molecular weight determinations. The sample is electrophoresed perpendicular to the urea gradient, so that the molecules moving across the gel are exposed to different urea concentrations. At positions where the urea concentration is high enough to promote unfolding, the mobility of the protein decreases because of the greater hydrodynamic volume of the unfolded form. The electrophoretic pattern thus obtained can be interpreted directly as an unfolding curve. Urea gradient gel patterns depend upon several parameters such as net charge, hydrodynamic volume, and conformational stability; therefore, it can be particularly useful for comparing different forms of a protein. Alternatively, a non-denaturing system is employed for analysis of biomolecules in their native state. Under denaturing conditions, the electrophoretic mobility proteins do not reflect their conformation, which exists at physiological pH. Several other properties that occur at normal pH may also be altered, such as protein–protein interactions or aggregation of proteins. Therefore, when the proteins are required to be electrophoresed on the basis of their native charge and hydrodynamic size, they are run under native gels in the absence of any denaturant. Another advantage with native gels is that it is possible to recover proteins in their native state after the separation, retaining their biological activity.
7.6
GEL ELECTROPHORESIS OF PROTEINS
7.6.1
SDS PAGE
Sodium dodecyl sulfate polyacrylamide gel electrophoresis is used for separating proteins primarily on the basis of molecular weight which determines its migration rate through the gel. The sample is treated with an anionic
200 Bioanalytical Techniques detergent SDS and b-ME which are included in the sample buffer for complete denaturation of protein molecules into linear polypeptides. SDS is also included in the gel and reservoir buffer for maintaining denaturing condition throughout the electrophoresis. SDS PAGE can be carried out in either a continuous or a discontinuous gel system. Continuous system involves running the protein sample in a single gel with a uniform pH throughout the path. It consists of the same buffer ions in the sample, gel, and electrode vessel reservoirs. It is less in use nowadays and has been largely replaced by a discontinuous gel system. Discontinuous or multiphasic buffer system employs variable buffer composition and pH in the gel and electrode reservoirs. The gels are also used with two different concentrations: the resolving or separating gel having higher concentration and stacking gel having lower concentration. The pore size of the stacking gel is larger than the separating gel and allows free mobility to the proteins. The proteins are stacked in concentrated bands, few microns to a millimetre thick, in the stacking gel before entering the separating gel. This is advantageous because it allows concentration of even dilute samples of proteins applied in large quantity to the gels, that too giving fairly good resolution. The proteins treated with SDS carry negative charge on their surface, therefore, they are resolved in an anionic system in which the electrodes are arranged in such a way that the upper bath acts as cathode while the lower acts as anode, and anionic proteins are allowed to flow towards the anode. The reservoir buffer (composed of Tris-glycine) has pH 8.3, the buffer in the sample and stacking gel (Tris–HCl) has pH 6.8, and the buffer in the resolving gel (Tris–HCl) has pH 8.8. The ions present in the buffer differ in the magnitude of their charge. Some ions have greater charge magnitude than the proteins, while the others have a lesser charge magnitude than the proteins. The stacking gel buffer and reservoir buffer yield Cl– and glycinate ions. Glycinate ions in the stacking gel are poorly ionized at pH 6.8 having lesser mobility and so are called ‘trailing ions’ while the Cl– ions having higher mobility are called ‘leading ions’. As soon as voltage is applied, Cl– migrates ahead of glycinate forming a voltage gradient which in turn accelerates the glycinate ions. While migrating, these ions form a sharp boundary with a high voltage gradient behind it. Since mobility of proteins is intermediate between the two ions, this moving boundary tends to sweep the proteins and resultantly they become stacked. See Table 7.2. As the anions enter the resolving gel, they encounter high pH (8.8). At this pH glycine is rapidly ionized and its effective mobility increases. The glycinate ions overtake the proteins and ultimately establish a uniform linear voltage gradient within the gel. Also on entering the resolving gel, proteins encounter lesser pore size posing molecular sieving effect due to which their
Electrophoresis
201
Table 7.2 Buffer and gel compositions Buffer
Composition
pH
Stacking gel buffer
125 mM Tris–HCl 0.1 per cent SDS
6.8
Separating gel buffer
375 mM Tris–HCl 0.1 per cent SDS
8.8
Electrophoresis buffer
25 mM Tris 250 mM Glycine 0.1 per cent SDS
8.3
Protein gel sample loading buffer
50 mM Tris–HCl 2 per cent SDS 10 per cent Glycerol 1 per cent b-ME 12.5 mM EDTA 0.02 per cent Bromophenol blue
6.8
mobility is retarded and they are unstacked. The proteins then sort themselves within this gradient according to their charge and size. The gradient in pore size as well as pH of gel running buffer observed inside the resolving gel, causes significant sharpening of protein bands during migration, enhancing the resolution. Proteins migrating under electric field are arrested at different positions in the gel and segregate into discrete bands corresponding to the rate of migration depending upon their size. Smaller molecules run faster; therefore, their bands appear away from the loading point and vice versa. Under proper conditions, proteins differing in their size as little as 2 kDa can also be separated and distinctively visualized in the form of bands. Acrylamide gels are usually prepared using glass plates and spacers. As stated above, standard protein gels are typically composed of two layers: stacking gel and resolving gel. Resolving gel is poured at the bottom and stacking gel at the top of it, in vertical position. To provide a smooth surface and interface between the two types of gels, isopropanol is placed over the resolving gel while it polymerizes. After complete polymerization the isopropanol is poured off, the top of the gel is rinsed with deionized water, and any excess water is soaked up carefully. Then the stacking gel is poured over the resolving gel and simultaneously a comb is inserted to form the sample wells. Protein mixture to be electrophoresed should be concentrated to about 40–60 µg/ml and should be free from insoluble impurities which might lead to streaking during electrophoresis. Sample buffer consists of 1 M Tris–HCl (pH 6.8), glycerol, SDS, b-ME, and 1 per cent bromophenol blue. Tris is used for minimizing cyanate modifications of proteins and glycerol increases the density of the sample so that it can sink into the sample well. Sample is boiled with SDS and b-ME for complete denaturation and attainment of uniform
202 Bioanalytical Techniques negative charge. Bromophenol blue, an anionic dye, is used as a tracking dye for monitoring the progress of protein molecules during electrophoresis. The samples are run into the stacking gel at about 10–15 mA for 10–15 min, and after they enter the resolving gel, current can be increased to 25–30 mA until completion, which takes approximately 3–4 h on a 10 cm gel. After the gel has run completely, the stacking gel is removed. Proteins are usually stained in Coomassie brilliant blue R-250 (CBB) which is also an anionic dye, prepared in methanol acidified with glacial acetic acid. The gel is kept in ample amount of staining solution for 3–4 h on shaker and is destained using acidified methanol so that the background gel becomes clear and proteins bands can be visualized (Figure 7.12). The time period for staining and destaining vary depending on the gel concentration, thickness, and protein concentration. Usually a prestained molecular weight marker is run parallel to the sample proteins that help to determine the molecular weight of proteins by comparing the relative distance travelled by them.
DETERMINATION OF MOLECULAR WEIGHT OF A PROTEIN Protein sample is electrophoresed by SDS PAGE with a set of MW standards. As shown in Figure 7.13 (a) and (b), the relative mobility Rf with reference to a marker protein and a tracking dye is calculated as given below: Rf =
Distance (in mm) migrated by protein 39 ___________________________________ ___ = = 0.67 58 Distance (in mm) migrated by dye
Figure 7.12
Polyacrylamide gel bands after SDS PAGE
stained
with
CBB
showing
separated
protein
Electrophoresis
203
These values are plotted to give a calibration curve of log MW versus relative migration distance (Rf) which is interpolated to calculate the molecular weight of the proteins of unknown size. This relationship is linear for any gel concentration over a limited range of molecular weight. The linear curve can be described by the following relation: y = mx + b
7.9
where y is the log MW, m is the slope, x is the Rf, and b is the y-intercept. The accuracy of the calculated MW depends on the linearity of the relationship, represented by the value of r 2. The closer the r 2 value is to 1.0, the better the fit of the data points to a line. From Figure 7.13(b) y = – 1.9944x + 2.7824 Since Rf value of unknown protein (x) = 0.67
Figure 7.13
Determination of (a) R f value of unknown protein from calibration curve
protein
bands
and
(b)
MW
of
204 Bioanalytical Techniques and, y = log MW Therefore MW = 10y – 10
–1.9944 (0.67) + 2.7824
MW = 28.1 kDa
7.6.2
Isoelectric Focusing
Isoelectric focusing of proteins is based on their amphoteric property. Charge on a protein molecule is due to the ionization of terminal a-amino or a-carboxyl groups, as well as the nature and number of ionizable alkyl (R) groups. The net charges of proteins are determined by the pH of their local environment and owing to their amphoteric nature, proteins exhibit characteristic isoelectric pH or pI which is a steady state position at which the net charges of proteins become zero (Figure 7.14a). All proteins exist as positively charged species below their pI and negatively charged species above their pI. This property of the protein forms the very basis of separation under IEF. During IEF, the proteins are allowed to move through a linear pH gradient due to which their native charges change according to the pH of their surrounding environment. The positively charged proteins move towards cathode becoming progressively less positive and negatively charged proteins move towards anode becoming less negative until they reach their pI. At pI, proteins bear net zero charge and do not move anymore in the gel. In this process, if a protein diffuses away from its
Figure 7.14
A protein’s (a) pI and (b) IEF
Electrophoresis
205
pI, it would either gain or lose a charge and migrate back to its pI value. This is called “focusing” effect due to which the proteins separate and concentrate into bands by virtue of their characteristic pI values (Figure 7.14b). This results in a high sensitivity for detection so that even small charge differences can also be differentiated. The resolution of proteins is determined by the slope of the pH gradient and the electric field strength. In order to obtain a good resolution, the IEF is generally performed under high voltage (around 1000 V) for a constant number of volt-hours (V h), an integral of the volts applied over the time. Isoelectric points of most proteins fall in the range of 3–12 pH and can be estimated with a calibration curve using marker proteins. The key feature of IEF is an appropriate pH gradient. There are two methods to create pH gradient: carrier ampholytes and immobilized pH gradient (IPG) gels.
7.6.2.1
Carrier ampholytes
These are mixtures of small amphoteric buffer molecules with high buffering capacity near their pI. Their molecular weight ranges from 400 Da to 1000 Da and their amphoteric property is due to the presence of multiple amino and carboxylic groups. These have high conductivity and form a pH gradient with a spectrum of pIs between 3 and 10 under the influence of the electric field. When the voltage is applied across a carrier ampholyte mixture, the positively charged carrier ampholytes move towards the cathode in the decreasing order of their pI values, while the negatively charged carrier ampholytes move towards the anode in the increasing order of their pI values. In this way all the charged molecules align themselves in an orderly fashion according to their pIs between the extremes, and buffer their environment to the corresponding pH resulting in a continuous pH gradient. The carrier ampholytes are usually used at about 2 per cent concentration and are added directly into the IEF gels. Higher concentration (more than 3 per cent) ampholytes are difficult to remove from gels and interfere with staining. The carrier ampholytes added to the gels give rise to the problem of cathodic drift, that is, the movement of the pH gradient off the basic part of the IEF gel with time. With cathodic drift, the pH gradient gradually drifts off the basic side of the gel, forming a plateau in the centre of the pH gradient. The cathodic drift usually occurs due to long focusing times and can be controlled by determining the optimum time of focusing in V h. Carrier ampholytes have been used extensively in past years, but now have widely been replaced by immobilized pH gradients due to several limitations. These are mixed polymers that are not well characterized due to compositional variations in manufacturing. The pH gradients made by a carrier ampholyte
206 Bioanalytical Techniques are often unstable and show cathodic drift over time which adversely affects the reproducibility of results.
7.6.2.2
Immobilized pH gradients
Alternatively, nowadays immobilized pH gradients (IPGs) are used for eliminating the problems of gradient instability and poor sample loading capacity associated with carrier ampholyte pH gradient See. Figure 7.15. Immobilized pH gradients are fixed gradients made of acrylamido buffers, commercially termed Immobilines which are weak acids or bases. Their buffering property is due to carboxylic and tertiary amino groups. These Immobilines are covalently incorporated into the polyacrylamide gel at the time it is cast to create an immobilized pH gradient. Since the separation under IEF occurs only according to the charge, the gel matrix must contain large pore sizes. For this purpose, polyacrylamide gels made of 5 per cent T and 3 per cent C giving a pore diameter of 5.3 nm are commonly used matrices. Continuous change in the ratio of these Immobilines enable us to obtain all types of gradients from wide to very narrow range. A linear gradient ranging from pH 3 to pH 10 enables the study of total protein distribution and estimation of protein’s pI. For studying the protein pattern in more detail with high resolution, narrow pH gradients are available giving smaller ranges, such as pH 3.5–4.5, 4–5, 4.5–5.5, 5–6, etc. Figure 7.15 gives the general structure of immobiline reagents. When incorporated into the gel, the acrylamido buffers provide even and controlled buffering capacity and a uniformly low conductivity throughout the gradient. The pH gradient is cast into the polyacrylamide gel, which is supported by a plastic backing. The cathodic drift is eliminated due to covalent bonding and because the pH gradient is fixed during the gel casting step, rather than formed during the first part of electrophoresis, as with carrier ampholyte gels. IPG gels are much more reproducible than carrier ampholyte gels, and can focus greater amount of protein than carrier ampholyte gels (up to 5 mg or more). Using IPG gels, the electrophoresis can be carried out at much higher voltage potentials (up to 10,000 V) and for much longer V h since the pH gradients are fixed and remain stable. These give resolutions as high
Figure 7.15
General structure of immobiline reagents
Electrophoresis
207
as 0.001 pH units using ultra-narrow pH gradients under high electric field strength. Commercial pre-cast IPG gels are also widely available for highly reproducible results. The IEF system comprises IPGphor strip holders that serve both as rehydration and as IEF chambers, and the IPGphor unit, which includes a high voltage power supply and built-in temperature control. IEF is usually performed using a horizontal flatbed electrophoresis unit attached with a Peltier cooling plate or ceramic cooling plate connected to a thermostatic circulator for efficient cooling and temperature control. Horizontal apparatus supports application of high voltages to achieve a high field strength giving sharply focused bands and the cooling system allows IEF to be performed under high voltage (3500 V) and maintaining this voltage for at least several thousand V h. It also eliminates excessive heat generated in the gel which can affect the sensitivity of the system. Parameters such as rehydration temperature and time duration, focusing temperature and maximum current, and the duration and voltage pattern are programmable. Commercially available apparatuses are versatile and fit a wide range of IPG strips. Commercial IPG gels are supplied as dried gel strips which require rehydration prior to electrophoresis. Sample can be applied either by including it in the rehydration solution (rehydration loading) or by applying it directly to the rehydrated IPG strip via sample cups, sample wells, or paper bridge. The mode of sample application is dependent on the sample composition and the IEF gel type. Usually rehydration loading is preferable as it is simple and allows more dilute samples to be loaded. It also allows larger quantities of proteins to be loaded and separated (up to 1 mg of sample per strip can be diluted or dissolved in rehydration solution for IEF). A typical rehydration solution generally contains urea for denaturation and solubilization of proteins, non-ionic or zwitterionic detergents for solubilizing hydrophobic proteins and minimizing protein aggregation, DTT as reducing agent to cleave disulfide bonds for complete unfolding of proteins, IPG Buffer or Pharmalytes which are carrier ampholyte mixtures used for enhancing protein solubility and producing more uniform conductivity across the pH gradient without disturbing IEF or affecting the shape of the gradient. Tracking dye (bromophenol blue) is included which allows the IEF progress to be monitored. Isoelectric focusing in the Multiphor II system is conducted at very high voltages (up to 3500 V) and very low currents (typically 3 (Figure 7.23). The capillary wall is now negatively charged and develops a double layer of cations attracted to it from the buffer. The inner layer is stationary, while the outer layer is pulled in the direction of charged cathode causing a powerful bulk flow which draws all the cations, anions, and neutrals towards the cathode end where the detector is placed. The electro-osmotic velocity of cations is highest so these reach the detector first, followed by neutrals, and then anions. The electro-osmotic velocity is given as neo = meoE
7.15
where meo is the electrophoretic mobility under EOF, representing a constant of proportionality between electro-osmotic velocity, and electric field strength. Electrophoretic mobility meo also depends upon the dielectric constant of the solution (e), the viscosity of the solution (h), and the zeta potential (z). µeo = e z/4ph
7.16
The zeta potential is determined by the nature and density of charge on capillary walls and nature of ions in the buffer. The EOF is a “plug flow” having an even and flat front which accounts for large number of theoretical plates in the capillary, counting as high as
Figure 7.23
Diagrammatic representation of EOF inside a capillary tube
Electrophoresis
221
50,000–500,000 under ideal conditions. The number of theoretical plates, or separation efficiency, in capillary electrophoresis is given by N = µ net V / 2D
7.17
where N is the number of theoretical plates, mnet is the net electrophoretic mobility of a molecule, D is the diffusion coefficient of the molecule, and V is the voltage applied. The number of theoretical plates and thereby the efficiency of the capillary electrophoresis are dependent upon the strength of the electric field and is limited only by the longitudinal diffusion. The basic instrumental set-up of capillary electrophoresis (Figure 7.24) consists of a high voltage power supply of about 10–30 kV, a capillary tube, two buffer reservoirs, two electrodes, a sample introduction system, a detector, and an output device for data acquisition. The capillary tube is usually 10–100 µm in diameter and 0.2–1 m in length, and is made up of fused silica (SiO2) which is used due to its transparency over a wide range of electromagnetic spectrum and high thermal conductance. It can also be easily manufactured into capillaries of few micrometres. The capillary tube dipped in the buffer reservoirs on each side contains the electrode and buffer or electrolytic solutions (anolyte and catholyte) which contributes to the EOF.
Figure 7.24
Instrumental set-up of a capillary electrophoresis
222 Bioanalytical Techniques In a capillary electrophoresis, extremely small amount of sample needs to be injected, typically ranging from picolitres to nanolitres. Specific amount of sample is injected by controlling either the injection voltage or the injection pressure. Two injection methods can be employed: hydrodynamic injection and electrokinetic injection. Hydrodynamic injection is accomplished by application of pressure or vacuum difference between the two ends of the capillary. This pressure difference makes the liquid to move into the capillary. Electrokinetic injection works when the capillary is placed into the catholyte on one end and anolyte on the other end. When voltage is applied, the EOF moves from tip to end of the capillary. This effect helps in dragging a representative sample into the capillary. After sample application, the electric field is induced which initiates the migration of ions through the capillary tube. For detection of the samples, many CE detectors such as absorbance detector and fluorescence detectors are used which have their own limits of detection (LOD). For example, UV has a LOD of 10 –13 –10 –12. Fluorescence detectors have LOD of 5 × 10 –17 and mass spectrometer has an LOD of 1 × 10 –17 . There is usually a small window near the cathodic end of the capillary which allows UV–VIS light to pass through the analyte and measure the absorbance. A photomultiplier tube is also connected to the cathode end which enables mass spectrum, providing information about mass to charge ratio of ionic species. The output of the detector is sent to the data output and handling device such as integrator or computer. The data is then displayed as an electropherogram, which reports the detector’s response as a function of time. Separated compounds can be observed as peaks with different migration times in an electropherogram (Figure 7.25). Capillary electrophoresis is a versatile and rapid technique enabling complex analyses of a variety of different molecules ranging from small inorganic ions to large nucleic acid fragments and proteins. It can be used for pharmaceutical analysis of large number of drugs and related substances and is particularly advantageous due to its high sensitivity and usage of very less amount of sample (in microlitres). Capillary electrophoresis is now being extensively utilized in clinical biochemistry, haematology, bacteriology, virologyl and immunology, and molecular and genetic analysis for high through-put applications. Dynamically coated capillaries have been used for separation of basic, acidic, and neutral drugs with high precision and selectivity. There are several variants of capillary electrophoresis method employed according to specific requirements. Capillary isoelectric focusing (CIEF) is a technique of choice for separation of amphoteric molecules such as proteins, peptides, amino acids, and pharmaceuticals, which can be readily adapted to
Electrophoresis
Figure 7.25
223
Electropherogram of a capillary electrophoresis
the capillary environment. Their separation is based on the difference in their pIs which guide their movement through a pH gradient inside the capillary. Affinity capillary electrophoresis (ACE) is also a versatile technique used for studying a variety of bimolecular non-covalent interactions forming complexes, such as protein–drug, protein–DNA, peptide–carbohydrate, peptide–peptide, DNA–dye, carbohydrate–drug, and antigen–antibody. Capillary isotachophoresis (CITP) employs discontinuous electric field causing difference in mobility of electrolyte inside the capillary. The sample is applied between the leading electrolyte which is low conducting and has a lower mobility as well as the trailing electrolyte which is high conducting and has higher mobility. This creates sharp boundaries between the molecules. Separation occurs on the basis of differences in the velocities of ions within the sample zones. Another form of CE is micellar electrokinetic capillary chromatography (MECC) which is capable of separating both neutral and charged molecules through the use of micelles in the separation buffer. Micelles are aggregates of amphiphilic monomers made of a hydrophilic head and a hydrophobic tail. Solutes are partitioned between micelles and the aqueous buffer as a result of various types of interactions such as hydrophobic, electrostatic, hydrogenbonding, etc. This partitioning forms the basis of differential migration of ionic species and their separation in MECC.
8 Spectroscopy I
8.1
ELECTROMAGNETIC RADIATIONS The continuum of the electromagnetic spectrum encompasses all kinds of radiation, from low energy radio waves and microwaves to very high energy X-rays and g rays (Figure 8.1). Visible region extending from 200 nm to 800 nm of wavelength, recognized by different colours as VIBGYOR, is the most familiar part of the spectrum. Electromagnetic waves travel in a defined direction with a constant speed that is the speed of light (c). Each wavelength corresponds to its frequency which is defined as the number of waves passing through a point each second. The energy associated with each electromagnetic radiation is proportional to its frequency. The relationship defining frequency (u) and energy (E) of an electromagnetic radiation is given as
Figure 8.1
Chart depicting the electromagnetic spectrum and various spectroscopic techniques for respective regions
226 Bioanalytical Techniques u = c/l
8.1
and E = hu
8.2
where h is the Planck’s constant = 6.62 × 10 –34 Js. Spectroscopic techniques are essentially based on the interactions of electromagnetic radiations with matter governed by the quantum laws. Matter is composed of atoms and molecules which can reflect, refract, transmit, or absorb electromagnetic radiation of specific energy. Atoms and molecules exist in unique energy state called ground state and possess distinct energy levels depending upon their chemical structures, due to which certain specific transitions are permissible on absorption of corresponding amount of energy. High-energy X-ray photons have sufficient energy to eject out an electron from an atom and cause ionization. The UV–visible radiations occurring in wavelength band of 400–700 nm possess energy below the ionization threshold of an atom and are absorbed by an orbital electron. These radiations are able to cause energetically favoured electron transitions usually from the singlet ground state S0 to the singlet excited state S1, because the energy difference between these orbitals corresponds to the energies of quanta in the UV–visible region. Electromagnetic radiation having longer wavelength such as infrared, microwaves, and radio waves possess energy insufficient to cause ionization or even electronic transitions from S0 to S1 state. Each electronic energy level in an atom is further split into a number of vibrational and rotational levels, where the electrons can make transitions on acquiring energy from these radiations (Figure 8.2). Thus, these radiations are able to bring about some bond oscillations and vibrations or introduce rotational changes in the atoms and molecules. All these transitions can be recorded with the help of a suitable spectroscopic method and can be used for generating lot of information leading to the qualitative and quantitative determination of compounds. For convenience of discussion, the topic ‘‘spectroscopic techniques’’ are divided in two chapters in this book. The present chapter deals with the spectroscopic techniques related to the phenomenon of absorption of radiations in short wavelength region bringing about ionizations or electronic transitions such as UV–visible spectroscopy, fluorescence spectroscopy, atomic absorption spectroscopy, X-ray spectroscopy, circular dichroism, and optical rotatory dispersion.
Spectroscopy I
Figure 8.2
227
Energy-level diagram depicting electronic transitions
8.1.1 Absorption and Emission Spectra When an atom or a molecule encounters an electromagnetic radiation, it acquires some of its energy, gets excited, and transits from its ground state to a higher energy state depending upon the energy contained by the incident photon. The specificity of transition dictates the absorption of characteristic radiation by the atom which gives it the required energy to jump from the ground state to a higher energy level equivalent to the energy difference denoted as DE between the two levels (Figure 8.3). Thus, if a continuous spectrum of radiation is supplied to an element, only the wavelengths corresponding to possible energy transitions within its atoms will be absorbed which will be seen as dark lines in the spectrum referred to as absorption spectrum. A plot of wavelengths can be obtained at which a given compound shows absorption maxima or “peaks”, which are highly useful for identification and characterization of compounds absorbing in a particular wavelength range.
228 Bioanalytical Techniques
Figure 8.3
Representation of transition of an electron from ground state to excited state
The higher energy state of an atom achieved after absorption of energy from the electromagnetic radiation is unstable and, therefore, the excited electron of the atom tends to relax back to its ground state by emitting discrete photons of characteristic energies which give a unique pattern of emission lines for that element. A plot of emission against wavelength for any given excitation wavelength is known as the emission spectrum. The uniqueness of the emission spectrum so obtained is due to the specific kind of transition occurring in the element and the difference in energy between levels involved in the transition. The intensity and width of the lines directly represent the number of photons emitted. In a graphical representation of the spectrum, the emission of a characteristic energy is shown as a peak in which the height and width of the peak represent its intensity (Figure 8.4).
8.2
ULTRAVIOLET AND VISIBLE LIGHT SPECTROSCOPY Ultraviolet and visible spectroscopy is the most common spectroscopic technique used in chemical and biological sciences for identification and quantification of large number of compounds. It is based on absorption of radiation by compounds in ultraviolet and visible spectral region extending from 200 nm to 800 nm in wavelength. Far-UV light having wavelength less than 200 nm is not frequently used for chemical analysis because it is strongly absorbed by air. The absorption of UV or visible radiation by a molecule leads to excitation of valence electron from the ground state to the excited state. The transition of electron to higher energy electronic level happens due to its transfer from bonding or non-bonding orbitals. The difference in energy between molecular bonding, non-bonding, and antibonding orbitals ranges corresponds to the energy of electromagnetic radiation in the UV region and the visible region of the spectrum. In fact this spectral range encompasses the majority of such electron transitions which make the basis for absorption of characteristic frequencies by inorganic and organic compounds, allowing for their identification and quantification.
Spectroscopy I
Figure 8.4
8.2.1
229
Illustration of (a) a continuous spectrum and (b) an absorption and emission line spectrum with their graphical representations
Basic Principles
Molecules absorb energy from UV or visible light to excite the electrons from lower energy bonding orbitals called highest occupied molecular orbital (HOMO) to higher antibonding molecular orbitals, called lowest unoccupied molecular orbital (LUMO). Sigma (s) bonding orbitals involved in single bonds
230 Bioanalytical Techniques are lower in energy than pi (p) bonding orbitals which occur in double or triple bonds. Non-bonding orbitals (n) having lone pair electrons are of highest energy. Absorption of energy can cause transition from one of these orbitals to an antibonding orbital that is s* or p* (Figure 8.5). The s Æ s* transitions require large energy and do not show the absorption maxima in typical UV–visible spectrum in the range of 200–800 nm. However, n Æ s* transitions are typically shown by saturated compounds containing atoms with lone pairs, on absorption of UV wavelength (150–250 nm). The most widely studied transitions by UV–visible spectroscopy are n Æ p* and p Æ p* transitions carried out by unsaturated compounds on absorption of radiation in the spectral range of 200–700 nm. These transitions are easily measurable through this technique. The radiation absorbing species in a compound is called “chromophore” which shows its characteristic absorption maxima. The absorption spectrum of a chromophore, however, is largely affected by its environment such as polarity and pH of the solvent dissolving the chromophore. With increasing solvent polarity, peaks resulting from n Æ p* transitions are shifted to shorter wavelengths due to higher degree of solvation of the lone pair which lowers the energy of the n orbital. This is called blue shift or hypsochromic shift. Similarly, decrease in the solvent polarity may cause a slight decrease in the energy difference between the excited and the unexcited states for n Æ p* transitions, causing a red shift or bathochromic shift. The principle behind the operation of UV–visible spectrophotometer is Beer–Lambert’s law which relates the concentration of the solution to the absorbance. As the radiation of a particular wavelength passes through a
Figure 8.5
Energy-level diagram depicting transition from bonding and non-bonding orbitals to antibonding orbitals
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solution, some of its intensity is absorbed while rest of it is transmitted (Figure 8.6). The percentage transmittance (T%) is given as T% = I / Io × 100
8.3
where I is the intensity of the transmitted light and Io is the intensity of the incident light. The relationship between absorbance (A) and transmittance (T) is given as A = log10 (1/T) = log Io /I
8.4
Lambert’s law says that the absorption of light is directly proportional to the thickness of the medium in the light path. A a l
8.5
where l is the thickness of the medium in the light path. According to Beer’s law, the absorption of light is directly proportional to the concentration of the analyte present in the medium. A a c
8.6
where c is the concentration of the analyte present in the medium. Combining the two, we get Beer–Lambert’s law. It defines the relationship between absorbance (A) and transmittance (T) as A a c × l
8.7
or A = e c l where e is the molar absorptivity or the molar extinction coefficient of the compound (dm3 mol–1 cm–1).
Figure 8.6
Transmittance of light through a solution of a defined path length
232 Bioanalytical Techniques It is a constant and characteristic of a given absorbing species in a particular solvent at a particular wavelength. Thus, Beer–Lambert’s law states that when a parallel beam of monochromatic light passes through a liquid medium, the amount of light that is absorbed is directly proportional to the concentration of the substance. The plots of transmittance and absorbance against the concentration are given in Figure 8.7.
8.2.2
UV–Visible Spectrophotometer
A UV–visible spectrophotometer essentially consists of a radiation source, a monochromator, a sample holder, and a detector. The radiation source is meant to give a stable, uninterrupted, high intensity output covering a wide range of wavelengths. Commonly used radiation sources are tungsten filament lamp, deuterium arc lamp, and xenon arc lamp. Light emitting diodes (LED) are also used frequently for the visible wavelength of light. Radiation source is followed by a monochromator which is used to disperse the radiation of a selected wavelength. Prisms and replica gratings are used extensively for this purpose. Glass prisms disperse light strongly over the visible region of the spectrum; however, these are not transparent to UV radiation with wavelength between 350 and 200 nm. Quartz and fused silica prisms are efficiently used over the range of UV spectrum. The sample is placed in cells called cuvettes. These are typically rectangular or cylindrical in shape, constructed with an internal width or diameter of 1 cm, respectively, to provide the unit path length to the radiation. Most common cuvettes are made of high grade fused silica or quartz since these are transparent
Figure 8.7
Graph plotted between (a) concentration of a solution and transmittance, and (b) concentration of a solution and absorbance
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throughout the UV, visible, and near infrared range. Use of glass and plastic cuvettes is limited to visible wavelengths. The transmitted light passing through the sample is collected and measured by the detector. Commonly used detector is typically a photomultiplier tube, a photodiode, or a charge coupled device (CCD) which proportionately converts the intensity of light transmitted from the sample into current. The signals are transmitted to a recorder which plots absorbance against wavelength and shows absorbance peaks at a particular wavelength (Figure 8.8). Two designs of spectrophotometer are commonly available: single-beam spectrophotometer and double-beam spectrophotometer. In a single-beam instrument (Figure 8.9), light transmitted through either sample or a blank is measured one at a time. The use of “blank” is a common practice in analytical methods based on absorption characteristics of a compound, to eliminate any possibility of error due to absorption of radiation by any species other than the one in question. This procedure compensates also for reflection, scattering, or absorption of light by the cuvette and the solvent. In a double-beam spectrophotometer (Figure 8.10), light is split into two beams before it reaches the sample through a beam splitter located between the exit slit of a monochromator and the cuvettes. After splitting, one beam passes through the reference and the other through the sample. Some doublebeam instruments have two detectors (photodiodes), and the sample beam and the reference beam are measured at the same time. In other instruments, the two beams pass through a beam chopper, which blocks one beam at a time.
Figure 8.8
A plot between absorbance and wavelength, showing absorbance peaks of two hypothetical compounds A and B
234 Bioanalytical Techniques
Figure 8.9
Schematic representation of a single-beam UV–visible spectrophotometer
Diffraction grating
Rotating disc Mirror Sample cell
Slit
Detector and Computer
Light source Mirror
Reference cell
Chart recorder
Figure 8.10
Schematic representation of a double-beam UV–visible spectrophotometer
The detector alternates between measuring the sample beam and the reference beam in synchronism with the chopper.
8.2.3 Applications of UV–Visible Spectrophotometer 8.2.3.1
Quantitative estimation of compounds
The compounds absorbing in UV–visible region can be quantified by using this technique. This determination is based on Beer–Lambert’s law and can be applied to any concentration range exhibiting linear relationship with absorbance at a particular wavelength. It is in fact the most popular technique
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for quantification of large number of organic and inorganic compounds due to its simplicity and efficiency. (i) Estimation of concentration using molar absorptivity If the molar absorptivity at a particular wavelength is known, the concentration of a compound in a solution can be determined by measuring the intensity of light that a sample absorbs. The path length of light is already fixed to unity by taking the cuvette diameter as 1 cm. For example: If A = 1.84 and e = 187000 dm3 mol–1 cm–1, then concentration of the compound may be calculated by the formula: A = e × l × c Substituting the values in the formula, we get or 1.84 = 187000 × 1 × c or c = 1.84 / 187000 = 9.83 × 10 –5 mol dm–3 (ii) Estimation of concentration by plotting a calibration curve If the molar absorptivity of the compound is not known, the unknown concentration can be determined using a standard curve of some known compound which is a plot of absorbance versus a set concentration values within a range, usually prepared by serial dilutions (Figure 8.11). Determination of unknown concentrations of proteins and nucleic acids is a good example since these molecules show absorption maxima at 280 nm and 260 nm, respectively.
Figure 8.11
A plot of absorbance versus concentration
236 Bioanalytical Techniques 8.3.2.2
Qualitative analysis
As stated earlier, most of the compounds absorbing in UV–visible region show very characteristic absorption at a particular wavelength, on the basis of which these can be identified. Although UV–visible spectra do not enable absolute identification of an unknown compound, the task can be accomplished by comparing the absorption spectrum of an unknown compound with the one of known compound called a reference spectrum. The absorption spectrum curves are prepared by recording the absorption at different wavelengths. The wavelength of maximum absorption (l max) depends on the presence of particular chromophores in a molecule; for example, nucleic acids absorb at 260 nm and proteins absorb at 280 nm. Some complex organic compounds show absorption at several maxima and each of them will have a characteristic shape and range indicating the presence of a particular functional group, some of which are listed in Table 8.1. Figure 8.12 shows the absorption spectrum of several plant pigments. Each of the compounds has its own peculiar absorption spectrum which helps in its identification. (i) Measurement of enzyme activity Enzyme activity can be easily, quickly, and conveniently calculated when the substrate or product is coloured or absorbs light in the UV–visible range. The rate of appearance or disappearance of light absorbing product or substrate can be measured with the help of a UV–visible spectrophotometer. For example, lactate dehydrogenase is an enzyme involved in the transfer of electrons from lactate to NAD +. The reaction is shown as follows: Lactate + NAD +
Pyruvate + NADH + H+
8.9
Table 8.1 List of some complex organic compounds and the wavelengths of their absorbance maxima Compound
l max (nm)
Adenine
260
Guanine
246
Cytosine
267
Uracil
259
Chlorophyll a
420, 663
Chlorophyll b
450, 650
Tyrosine
275
Phenylalanine
260
Proteins
280
Nucleic acids
260
Spectroscopy I
Figure 8.12
237
Absorption spectrum of plant pigments
In this reaction, the products are pyruvate, NADH, and a proton. NADH absorbs radiation in the ultraviolet range at 340 nm and its counterpart NAD + does not. Neither any of the substrate nor the product absorbs at 340 nm. Thus, the progress of the reaction in forward direction can be followed by measuring the increase in absorption at 340 nm in spectrophotometer. (ii) Determination of dissociation constant Dissociation constants of acids and bases, and their pKa value can be calculated if the ratio of [A– ] / [HA] is known at a particular pH and can be determined spectrophotometrically from the graph plotted between absorbance and wavelength at different pH values. (iii) HPLC detector The UV–visible spectrophotometer is commonly hyphenated to automated chromatographic systems like HPLC. The presence of an analyte gives a response which can be assumed to be proportional to the concentration. For more accurate results, the instrument’s response to the analyte in the unknown sample should be compared with the response to a standard. (iv) Determination of impurities Impurities in organic molecules can be determined with the help of additional peaks observed in the spectrum
238 Bioanalytical Techniques which can be compared with the absorbance of standard material at a specific wavelength.
8.3
FLUORESCENCE SPECTROSCOPY Fluorescence spectroscopy has become quite popular because of its acute sensitivity that enables it to identify even the slightest change in the structural and dynamic properties of biomolecules. Fluorescence measurements can provide information on numerous molecular processes extending its applications from chemical to biological research. Intracellular processes, such as biomolecular interactions, metabolic reactions, conformational changes, and intracellular localization can be detected and quantified more precisely and accurately on the basis of intrinsic or extrinsic fluorescence of compounds. Fluorescence spectroscopy is now a popular technique used extensively in biotechnology, medical diagnostics, DNA sequencing, forensics, and genetic analysis. Fluorescence occurs on excitation of a fluorophore when it absorbs an appropriate wavelength of light followed by the emission of a photon to return back to its ground state from an excited one (refer to Figure 3.13). Fluorophore, excited to higher vibrational level usually S1 or S2, tends to relax rapidly to the lowest vibrational level, at S1 through a process called internal conversion. Thereafter, it returns quickly to the ground state. This transition generally occurs within 10 –12 s or less. Fluorophores are usually organic compounds which emit light of specific wavelength after absorption, while the compound makes transition from a higher energy level to a lower energy level. The energy of this emitted light (fluorescence) depends on the chemical structure of the fluorophore as well as on its surrounding environment. The efficiency of a fluorophore to absorb the excitation light is given by its extinction coefficient denoted by E. The yield of emitted light or fluorescence, also called the quantum yield, is defined as the ratio of the quanta of light emitted to the quanta of light absorbed. In order to obtain maximum fluorescence intensity, a fluorophore is usually excited at the wavelength at the peak of the excitation curve, and the emission is detected at the peak wavelength of the emission curve by using appropriate filters. Fluorophores exhibit their own specific absorption and emission spectrum depending on the internal structure of the fluorescing molecule and binding chemistry (Figure 8.13). Commonly used fluorophores are fluorescein, rhodamine, acridine, and quinine (Figure 8.14).
Spectroscopy I
Figure 8.13
Excitation and emission spectra of fluorescein
Figure 8.14
Some common biochemical fluorophores
8.3.1
239
Basic Principles
Intrinsic or extrinsic fluorescence of compounds can be utilized for their quantitative determination since it is related to the concentration of the absorbing species by Beer–Lambert’s law, as stated here
240 Bioanalytical Techniques Io / I = e
–Ecl
8.10
or log10 Io / I = Ecl
8.11
where I is the intensity of transmitted light, Io is the intensity of incident light, E is the molecular extinction coefficient, c is the concentration in g M L –1, and l is the path length of sample. log10 Io /I denotes the absorbance or optical density of the sample. Direct measurement of fluorescence intensity can be used to determine the concentration of the analyte. However, as the concentration of fluorophore increases, deviations occur and the plot of emission against concentration becomes non-linear. The two characteristic features of fluorescence are lifetime and quantum yield (Q) of fluorescence. The lifetime of the fluorescence is usually 10 ns, defined as the average time the fluorophore spends in the excited state prior to returning to the ground state. The fluorescence quantum yield is defined as the ratio of the number of photons emitted to the number of photons absorbed. Compounds having high quantum yields (ª1) show brightest emissions; for example, fluorescein and rhodamines. Number of photons emitted Q = __________________________ Number of photons absorbed
8.12
The fluorescence emission spectrum also reveals information about the environment of the fluorophore and its interactions. Fluorescence typically occurs at longer wavelengths, that is, fluorescence photons are longer in wavelength than the excitation radiation and this is referred to as Stokes shift. The Stokes shift accounts for the sensitivity of fluorescence techniques because it allows emission photons to be detected against a low background, isolated from excitation photons. Stokes shift is mainly caused by the decay of fluorophore by rapid dissipation of energy during transition from excited state to the lowest vibrational level of S1. Since fluorophores finally decay to higher vibrational levels of S0, it leads to further loss of excitation energy by thermalization of the excess vibrational energy. Other causes of Stokes shifts are fluorophore–solvent interactions, complex formation, and other modes of energy transfer.
Spectroscopy I
8.3.1.1
241
Fluorescence quenching
Quenching refers to the decrease in intensity of fluorescence which occurs due to the presence of some other molecules in solution known as quencher. Quenching may occur either through energy transfer or through electron transfer from fluorophore to quencher. Attenuation of the incident light by the quencher or fluorophore itself may also lead to reduction in fluorescence intensity. Collisional quenching occurs during the excited state lifetime of a fluorophore, when the excited fluorophore collides with a quencher bringing about non-radiative transitions. Common collisional quenchers are O2, I–, Cs+, and acrylamide. Another usual form of quenching is static quenching in which quenchers may form stable complexes with some molecules of fluorophore at the ground state. Rest of the fluorophore molecules which remain unaffected are excited by the incident radiation and emit the fluorescence normally. This decreases the net intensity of fluorescence depending upon the concentration of the quencher.
8.3.1.2
Fluorescence spectrum
Most spectrofluorometers are able to record both excitation and emission spectra. The fluorescence spectral data represented as an emission spectrum is a plot of the fluorescence intensity versus wavelength (nm) or wavenumber (cm–1). Conversely, an excitation spectrum is a plot of excitation wavelengths measured for a single emission wavelength. Emission spectra are typically independent of the excitation wavelength. For most fluorophores the emission spectrum is typically a mirror image of the absorption spectrum of the S0 Æ S1 transition because same transitions are involved in both absorption and emission, and vibrational energy levels of S0 and S1 are similar. The symmetric nature of absorption and emission spectrum is depicted in Figure 8.15.
8.3.2
Fluorescence Spectrometer
Fluorescence spectrometers consist of a light source, monochromators, and high sensitivity detectors. The most common light source is xenon arc lamp that provides a continuous UV–visible light output of high intensity at all wavelengths. High-pressure mercury lamps and high-pressure mercury–xenon lamps are also used as high intensity sources. The light from an excitation source passes through an excitation filter or monochromator which selects a specific wavelength and eliminates the stray light. Most fluorescence spectrometers employ diffraction gratings usually holographic gratings as monochromators, having high dispersion efficiency and stray light rejection. The excitation
242 Bioanalytical Techniques
Figure 8.15
Absorption and emission spectrum of a hypothetical compound
light after passing through the monochromator strikes the sample so that the molecules in the sample fluoresce. The fluorescent light is emitted in all directions. Some of this fluorescent light passes through a second filter or monochromator and reaches a detector, which is usually placed at 90° so that the transmitted or reflected incident light does not reach the detector. Fluorescence spectrometers generally use photomultiplier tubes (PMTs) as detectors. The photocathode in PMT is sensitive to longer wavelength emissions and is capable of detecting individual photons. The output from the detector is displayed on a readout device or digital display. The instrumentation scheme of a fluorescence spectrometer is depicted in Figure 8.16.
8.3.3 Applications of Fluorescence Spectrophotometer Fluorescence spectroscopy is now an indispensable investigational tool in many areas of science, due to its extremely high sensitivity and selectivity.
Figure 8.16
Diagrammatic representation of a fluorescence spectrophotometer
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The measurements of fluorescence spectrum, lifetime, and polarization are powerful methods enabling the analytical implementations of this technique in chemical, biochemical, and medical research allowing real-time analysis of dynamic processes. Using fluorescence resonance energy transfer (FRET) it has become easier to monitor protein–ligand binding and ascertain protein domain structure. Fluorescence spectroscopy is a promising diagnostic technique with high sensitivity and specificity for microorganisms associated diseases diagnosis with the help of spectroscopic fingerprints. It can be used for the determination of drugs in low dose formulations and in studying their binding properties. Fluorescence spectroscopy is also used in environmental sciences for determination of inorganic ions and heavy metals, as well as for testing petroleum pollutants in marine environment.
8.4 ATOMIC ABSORPTION SPECTROMETRY Atomic absorption spectrometry (AAS) is a sensitive technique used for determining the trace concentration of elements like metals. It is based on the fact that atoms of different elements absorb particular wavelengths of light, get excited, and produce characteristic spectra. The wavelengths of the light transmitted or absorbed by the sample are measured by a detector and, following Beer–Lambert’s law, the concentration of element is determined from a calibration curve, obtained by using standards of known concentration. More than 70 different elements can be quantified by this technique in diverse samples including biological and non-biological material. Its common most application is measuring elemental pollutants like heavy metals in environmental samples such as soil, water, organic material. This technique is extremely useful in analysing traces of metal in biological fluids such as blood and urine, and any metallic contaminants in pharmaceuticals, food products, beverages, etc.
8.4.1
Basic Principles
Atoms in their “ground state” absorb light energy of a specific wavelength and make transitions to higher energy levels called the “excited state.” As the number of atoms in the light path increases, the amount of light absorbed also increases; this provides the basis of a quantitative determination of an element by atomic absorption spectrophotometer. This instrument measures the amount of light absorbed by the sample through a detector and compares it to the light that originally passed through the sample. A signal processor then transforms these changes into signals which appear as energy absorption peaks at discrete wavelengths. Every atom has its own distinct pattern of absorption of characteristic wavelengths due to the unique configuration of electrons in its
244 Bioanalytical Techniques outer shell. This enables the qualitative analysis of a sample. The use of special light sources and careful selection of wavelengths allow specific determination of individual elements.
8.4.2 Atomic Absorption Spectrometer Atomic absorption spectrophotometer is often used to measure very minute quantities of elements in samples (ppm or ppb). Even slight contamination during sample handling or sample preparation can lead to severe error. Usually vessels made of inert material such as perfluoroalkoxy polymers (PFA), silica, teflon are used for storage so that the element does not adsorb or adhere to the surface of vessel. Sample may require acid digestion, filtration, centrifugation, or may be used directly for atomization. The basic components of an atomic absorption instrument include the light source that emits the desired radiation, an atomizer in which atoms of the sample are produced, a monochromator for light dispersion, and a detector which measures the light intensity and amplifies the signal to be read on a display (Figure 8.17). The molecules of sample are vaporized and dissociated into free atoms through a process called atomization. Most commonly used atomizers are flame atomizers and electro-thermal atomizers. Flame atomizers consist of a burner, a nebulizer, and a spray chamber. Burners are operated with two flame gases–an oxidant/fuel combination, such as air–acetylene and nitrous oxide–acetylene. Air–acetylene flame gives temperature of about 2300°C and can be used for determining approximately 45 elements by atomic absorption. Monochromator
Detector
Hollow cathode lamp
Nebulizer
Figure 8.17
Sample
Schematic representation of atomic absorption spectrophotometer
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The nitrous oxide–acetylene flame gives a maximum temperature of about 2900°C and is used for determining elements which form refractory oxides. The sample solution is aspirated pneumatically by the nebulizer and converted into aerosol. This aerosol is mixed with flame gases in a spray chamber. Electrothermal atomization uses a hollow graphite tube accommodating about 10–25 µl of sample. It is heated sequentially up to 2500°C by passing electric current through the tube in a controlled manner, rendering all the material in vaporized and atomized state. The radiation source is a narrow-line source of high intensity that provides the atoms of a particular element with very specific wavelengths at which absorption occurs. The sources used for atomic absorption are the hollow cathode lamp (HCL) and the electrodeless discharge lamp (EDL). Of these, a hollow cathode lamp is the most common one. It consists of a cathode made of an element to be determined and an anode, both sealed in a glass tube filled with an inert gas such as neon or argon at low pressure (about 1 Nm–2 to 5 Nm–2). On applying voltage across the electrodes, the gaseous ions hit the cathode and eject metal atoms in the process called ‘sputtering’. These sputtered metal atoms are in excited states. As they drop back to their ground state, they emit a few characteristic radiations, which pass through a monochromator. A monochromator disperses the light and selects the specific wavelength of light that is, spectral line, which is absorbed by the sample, and excludes other wavelengths. The light selected by the monochromator is directed onto a detector, which produces an electrical signal proportional to the light intensity. Photomultiplier tubes are the most common detectors for AAS.
8.4.3
Calibration of Atomic Absorption Spectrometer
The elemental quantification in a sample can be done using a calibration curve (Figure 8.18). Since the instrument relies on Beer–Lambert’s law, a plot of concentration versus absorbance can be made using a concentration gradient of a standard within the stipulated range. However, for most elements, particularly
Figure 8.18
Calibration curve plotted between concentration and absorbance
246 Bioanalytical Techniques at high concentrations, the relationship between concentration and absorbance deviates from Beer’s law and is not linear. The sample solution is fed into the instrument and the unknown concentration of the element is then displayed on the calibration curve. Samples may need dilution so that they fall in the linear range of the calibration curve.
8.5
X-RAY SPECTROSCOPY Highly penetrable electromagnetic radiations, which were later recognized as X-rays, were discovered by the German physicist Wilhelm Conrad Röntgen. He was awarded the Nobel Prize for his discovery in 1901. X-rays are very highenergy radiations, characterized by the wavelength range of 0.01–10 nanometres corresponding to the frequencies of 3 × 1016 Hz to 3 × 1019 Hz. The X-rays with wavelengths longer than 0.1 nm are called soft X-rays, while those with shorter wavelengths are known as hard X-rays. High sensitivity and flexibility of the technique makes it ideal for routine as well as non-routine specialized analyses. Most common use of X-rays has been in medical imaging. Besides, X-ray spectroscopy is now used in almost every field such as in elemental identification and quantification in medicine and pharmaceuticals, agriculture, forensics and material science, determination of alloy composition in manufacturing processes, for investigating toxic metals and other pollutants in soil and water in environmental sciences and much more. This technique is especially used in analyses of mineral ores such as calcite, dolomite, gypsum, manganese, and iron ores. XRF being a non-destructive technique is especially useful in analysis of antique items, paintings, ancient sculptures, archaeological artefacts, etc.
8.5.1
Production of X-Rays
X-rays are produced by interactions between electrons and atoms or between charged particles and an electromagnetic field. Bremsstrahlung and synchrotron radiation are used in X-ray spectroscopy as sources of continuous X-ray radiation. In the Bremsstrahlung process, a charged particle such as an electron is decelerated under an electric or magnetic field by another charged particle such as an atom. This changes the magnitude and direction of its velocity and results in loss of kinetic energy. By the law of conservation, this energy is converted into an X-ray photon. Bremsstrahlung is, therefore, also known as a “braking radiation” or “deceleration radiation”, and is a common source of a continuous spectrum of X-rays. A synchrotron consists of an electron storage ring as a radiation source, in which electrons are accelerated in circular orbits by strong electric and magnetic fields generating extremely bright X-rays. Synchrotron is known for its high
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spectral brilliance which is a measure of the photon flux per second, per unit cross-sectional area of the radiation source, per unit angle of the radiation cone, and per unit spectral bandwidth. Brilliance in spectroscopy is directly related to the spectral resolution which is achieved when appropriate radiation flux from the source is able to pass through the monochromators and slits, and is incident on the dispersing sample at uniform angle. Brilliance is especially important in crystallography for obtaining good resolution of closely spaced unit cells behaving as diffraction spots.
8.5.2
Interaction of X-Rays with Matter
Similar to the other electromagnetic waves, X-rays also interact with matter in a variety of ways depending on the energy of the electromagnetic radiation as well as the physical and chemical composition of the material. On encountering X-rays, matter may show its inelastic scattering called Compton scattering, elastic scattering also known as Rayleigh scattering, or may cause its photoabsortion. The phenomenon of photoabsorption predominates in the region of soft X-rays and causes photoelectric effect in which an X-ray is absorbed by the atom by transferring all of its energy to an innermost electron ionizing the atom and producing a photoelectron that is likely to ionize more atoms in its path. The electron deficiency can be filled by dropping of an outer electron of that atom into the vacant shell. During this process, the atom de-energizes itself by emitting a characteristic X-ray photon or an Auger electron (Figure 8.19). These effects can be used for elemental detection through X-ray spectroscopy or Auger electron spectroscopy. This electron is ejected with a characteristic kinetic energy which depends on the binding energy (Ebinding) of the electron in its orbital. hv = Ebinding +
8.5.3
1 __ mv2 2
8.13
X-Ray Absorption Spectroscopy
An X-ray absorption spectrum (XAS) is meant to study the intensities of incident and the transmitted X-ray which follows Beer’s law. If a monochromatic beam of X-ray of intensity I0 passes through a sample of thickness x, the intensity of the transmitted beam is I, then according to Beer’s law it will be given as ln (I0 /I) = µ x
8.14
where µ is the linear absorption coefficient, which depends on the atomic number and the density r of the material.
248 Bioanalytical Techniques
Figure 8.19
Emission of X-ray photon
In X-ray absorption spectrum, such absorption is seen to increase abruptly at certain energies and gives rise to an absorption edge where transmitted X-ray intensity drops drastically. The edge occurs when the energy of the incident photon is sufficient to cause excitation of a core electron of the absorbing atom and produces a continuum photoelectron. The edge energies, therefore, correspond to the binding energies of electrons in the K, L, M,..., shell of the absorbing atom (Figure 8.20). The absorption edges are named in the order of increasing energy, K, L I, L II, L III, M I,..., corresponding to the excitation of an electron from the 1s (2S½), 2s (2S½), 2p (2P½), 2p (2P3/2), 3s (2S½),…, orbitals and so on. An absorption edge gives vital information enabling elemental identification, however keen observation of the regions in immediate vicinity or a little farther can yield huge amount of information for structural analysis of the material. XAS, therefore, is not limited to the study of merely evident absorption edges but also refers to the measurement of X-ray absorption cross-section in the vicinity of one or more absorbing edges, more appropriately now referred to as X-ray absorption fine structure (XAFS). X-ray absorption fine structure can be studied at the “near edge” region that is, the structure in the vicinity of the edge, and is referred to as X-ray
Spectroscopy I
Figure 8.20
249
Schematic illustration showing (a) energy-level diagram showing major transitions and (b) X-ray absorption spectrum showing K, L, and M edges
absorption near edge structure (XANES) or near edge X-ray absorption fine structure (NEXAFS). In this region, transition of a core electron occurs to a close non-bound level having very low energy difference (Figure 8.21). Such transitions have high probability and the ejected photoelectron can be scattered by many other atoms in vicinity, prior to returning to the absorbing atom. Since the photoelectron produced in such a transition has a very low kinetic energy and consequently a long mean-free path, it experiences strong multiple scattering. X-ray absorption can also be studied in extended region, referred to as extended X-ray absorption fine structure (EXAFS), characterized by the photoelectrons having high kinetic energy and predominantly experiencing single scattering by the atom in close vicinity (Figure 8.21). EXAFS spectra are recorded typically in a range of 500–1000 eV from the absorption edge. In the process of photoelectron ejection, considering the wave nature of moving electrons, the back-scattered electron waves may interfere with the ejected electron waves. If the two waves are in phase at certain energy, it results into constructive interference observed as absorption maximum. If the two waves interfere destructively at a certain absorber–scatterer distance, it results into minimum absorbance. The interference patterns are modulated as absorption coefficient and reflect as oscillations in the EXAFS spectra.
8.5.4
X-Ray Fluorescence Spectroscopy
During photoelectric absorption, the X-ray photon transfers its energy to one of the lower shell (K or L) electrons of an atom, which is consequently ejected out
250 Bioanalytical Techniques
Figure 8.21
Schematic illustration of an X-ray absorption spectrum, showing XANES and EXAFS
of the atom. There creates an electron vacancy in the inner shell, which can be filled by an electron dropping from a higher energy shell. This is accompanied by the emission of an X-ray photon termed as secondary X-ray photon. The energy of this photon characteristically corresponds to the energy difference between the two shells and is called X-ray fluorescence. The secondary X-rays produced in this process are characteristic of a specific element, and can be used to identify the elements present in a sample by analysing the energy of the secondary X-rays. Photoelectric effect gives rise to an emission spectrum of X-rays at a few discrete frequencies referred to as the spectral lines. These characteristic spectral lines are named on the basis of orbitals involved in ion transitions, such as L shell to K shell transition is called K a , M to K shell transition is called K b , and so on (Figure 8.22). The intensity of each spectral line depends on the concentration and the atomic number of the fluorescing element in the sample, and provides vital information about the elemental composition of an unknown sample, which can be used for its qualitative and quantitative measurements.
8.5.5
X-Ray Diffraction
The X-ray diffraction (XRD) technique is based on the diffraction of an incident monochromatic beam of X-rays by the atomic planes of a crystalline material. A crystal is said to be made up of a periodic arrangement of several unit cells into a lattice. Each unit cell further consists of one or more atoms in a fixed arrangement. There are many atomic planes in a crystal, each with a definite spacing. Crystalline materials scatter X-rays producing interesting patterns through constructive and destructive interference. This was first discovered in the crystals of NaCl, KCl, and ZnS by W. L. Bragg and W. H. Bragg for which they were awarded Nobel Prize in 1915. Bragg’s law postulates
Spectroscopy I
Figure 8.22
251
Diagrammatic representation of transitions giving rise to X-ray emission lines
that X-rays having wavelength equal to an integer multiple of distances between these atomic planes are diffracted such that the diffraction angle is equal to the angle of incidence (Figure 8.23). 2d sinq = n l where d q n l
is is is is
8.15
the spacing between the atomic planes, the angle between the incident ray and the scattering planes, an integer, and the wavelength of the incident wave.
The directions of probable diffractions depend on the size and shape of the scatters, and the intensities of the diffracted waves depend on the type and arrangement of atoms in the crystal structure.
Figure 8.23
Depiction of Bragg’s diffraction
252 Bioanalytical Techniques X-ray diffraction is recorded by X-ray diffractometer (Figure 8.24), essentially having three basic elements: an X-ray tube, a sample holder, and an X-ray detector. The X-ray tube is a cathode ray tube which generates X-rays by heating a filament to produce electrons. These electrons are accelerated to bombard a target by applying voltage. Electrons of sufficiently high energy dislodge inner shell electrons of the target material and X-ray spectra of characteristic wavelengths are produced. Commonly used target material is copper that gives Ka radiation of wavelength 1.541Å. The X-rays are then passed through a filter to produce monochromatic beam of radiation which is collimated and directed onto the sample. The sample is mounted on a goniometer, an instrument which rotates the object to a precise angular position. Thus, the sample is maintained at an angle q with respect to the collimated X-ray beam. The X-ray detector is mounted on an arm in the X-ray diffractometer which rotates at an angle of 2q and collects the diffracted X-rays. The interaction of X-rays with the sample oriented in appropriate geometry satisfying Bragg’s law produces constructive interference and gives a peak in intensity. This X-ray signal is recorded by the detector, and is processed and counted on output device. X-ray diffraction analysis yields lot of structural information about the crystal lattice. It is a non-destructive analytical technique through which the unique fingerprint of Bragg diffractions can be obtained leading to determination of type and arrangement of atoms and identification of phases within a sample.
Figure 8.24
Schematic representation of X-ray diffractometer
Spectroscopy I
253
Since a crystal exhibits a set of unique d-spacings, conversion of the diffraction peaks to d-spacings leads to identification of the compound by comparison of d-spacings with standard reference patterns.
8.6
CIRCULAR DICHROISM AND OPTICAL ROTATORY DISPERSION Circular dichroism (CD) and optical rotatory dispersion (ORD) are closely related phenomena exhibited by optically active molecules or chiral molecules. These molecules exist as enantiomers which are mirror image isomers and not superimposable. Being asymmetrical, these molecules show difference in the absorption of right- and left-handed circularly polarized light which makes the basis of CD spectroscopy. A plane polarized light is said to have two vector components, electric (E) and magnetic (M), propagating perpendicular to each other (Figure 8.25). When the direction of the E-component is confined to a single plane perpendicular to the direction of propagation while its magnitude oscillates, it is called linearly polarized light. If the magnitude of the oscillation is held constant and the direction of the electric field vector rotates about its propagation direction, it is called a circularly polarized light. The electric vector of a circularly polarized light propagates in the form of a helix along the direction of propagation. If the electric vector rotates counter-clockwise, it is called left circularly polarized light (LCP), and if the electric vector rotates clockwise, it is called right circularly polarized light (RCP).
Figure 8.25
Plane polarized light travelling along x-axis, with its electrical vector along y-axis and magnetic vector along z-axis
254 Bioanalytical Techniques 8.6.1
Basic Principles
Circularly polarized light is generated by passing a plane polarized light (propagating in z-axis) through an optical element, a birefringent plate, which splits the plane polarized light into two plane polarized beams oscillating along different axes (x and y). This plate creates the phase difference of a quarter wavelength by retarding one of the beams. The two beams which are now 90º out of phase are added together resulting in a circularly polarized light propagating one direction (Figure 8.26). As the circularly polarized light passes through an optically active sample, the right circularly polarized light and left circularly polarized light are absorbed to different extents at a wavelength due to differences in their extinction coefficients. Therefore, CD is observed only for wavelengths where the substance absorbs light. The direction of the E-vector now traces an ellipse instead of a circle also causing the rotation of the major axis of the ellipse due to differences in refractive indices. This is referred to as CD. Circular dichroic absorbance can be given as DA = AL – AR
8.16
where AL is the absorbance of left circularly polarized light and AR is the absorbance of right circularly polarized light. The absorbance of both right circularly polarized light and left circularly polarized light follows the Beer–Lambert’s law, that is, I = I0e –A ln10 where I is the transmitted light intensity and
Figure 8.26
A right-handed circularly polarized light
8.17
Spectroscopy I
255
Io is the incident light intensity. Molar CD at a given wavelength l is given as the difference in molar extinction coefficients DA 8.18 De = eL – DR = _____ C × l where De is the molar CD or molar differential dichroic absorptivity expressed in litre mole−1 cm−1, eL is the extinction coefficient/molar absorptivity of left circularly polarized light, eR is the extinction coefficient/molar absorptivity of right circularly polarized light, c is the concentration of the test solution in mole litre−1, and l is the optical path of the cell in centimetres. Circular dichroism is also expressed in ellipticity (degrees cm2 decimole−1) and can be shown diagrammatically as Figure 8.27. 2.303 qr = ______ × (AL – AR ) × [rad] 4 or 2.303 180 qr = ______ × (AL – AR ) × ____ p × [deg] = 33.0047 DA 4 The molar ellipticity [qM] also expressed in degrees cm2 decimole−1 is given as
Figure 8.27
Elliptically polarized light
256 Bioanalytical Techniques q × M [qM] = _________ c × l × 10
8.19
where q is the value of ellipticity given by the instrument, M is the relative molecular mass of the substance to be examined, c is the concentration of the solution to be examined in g/ml, and l is the optical path of the cell in centimetres. Molar ellipticity is often used in the analysis of proteins and nucleic acids, and the molar concentration is expressed in terms of monomeric residue. If De or DA or ellipticity q is plotted against wavelength l, a CD spectrum can be obtained. Optical rotatory dispersion measures a change in the optical rotation with respect to the change in wavelength. In an optically active material the left and right circularly polarized components propagate with different speeds because the index of refraction (n) for the two components is different, that is, nL π n R
8.20
where nL is the index of refraction of left circularly polarized component and nR is the index of refraction of right circularly polarized component. This effect is called circular birefringence (CB) which results in a phase shift (d) between the two components, proportional to the difference in refractive index. The relation can be represented by the following expression: 2p ___ (nL –nR ) 8.21 l On superposition of the two components after passing through the distance (d) in the optically active medium, the phase shift results in a permanent rotation of the long axis (major axis) of the elliptically polarized light by an angle a. The angle of rotation a¢ in radians is calculated as follows d =
d 8.22 a¢ = __ = p /l (nL – nR )d 2 Optical rotatory dispersion is a dispersive technique which does not necessarily require the sample to absorb the wavelength at which it is being measured.
8.6.2
Circular Dichroism Spectrometer
A circular dichroism spectrometer (Figure 8.28) consists of a light source, a monochromator, an electro-optic modulator, and a detector. The light source
Spectroscopy I
Linearly Unpolarized Circularly polarized light polarized Light Polarizer light Electro-optic light source Monochromator modulator
Figure 8.28
Optically active sample
257
Detector
Schematic representation of a CD spectrometer
is usually a xenon lamp giving a continuous radiation of wavelength ranging 170–800 nm within which most dichrographs are obtained. The light passes from the source to a monochromator having quartz prisms which allow specific wavelength to pass through. The monochromatic light now passes through a polarizer where it is converted into a circularly polarized light. An electrooptic modulator is installed in the light path which allows either the left- or right-handed component of a circularly polarized light to pass to the sample. The chiral molecules in the sample cause the electric vectors of left- or righthanded components (denoted as EL and ER respectively) to oscillate elliptically according to their molar absorptivity. The detector, commonly a photomultiplier tube, generates the voltage proportional to the ellipticity of resultant beam.
8.6.3 Applications of Circular Dichroism Spectroscopy Circular dichroism spectroscopy is mainly concerned with the properties of chirality in molecules. Biological macromolecules like nucleic acids and proteins exhibit high degree of chirality due to their structural conformations and are, therefore, widely studied using this technique. Since CD spectra of protein and DNA are greatly influenced by their three-dimensional structures, they are mainly used for identifying dynamic changes in their structure induced by temperature, pH, denaturants, etc. It also helps in estimating secondary structure of polymers especially proteins and polypeptides. Further CD analysis allows determination of folding and unfolding of proteins, protein–ligand interactions, protein–nucleic acid, and protein–protein interactions. These are also widely used for studying denaturation, renaturation, supercoiling, DNA ligand interactions, and effects of mutations. CD spectra of some common conformations of proteins and nucleic acids are illustrated in Figure 8.29.
258 Bioanalytical Techniques
Figure 8.29
Graphical illustration of CD spectra of (a) secondary conformations of proteins and (b) A, B, and Z forms of DNA
9 Spectroscopy II
9.1
INTRODUCTION The longer wavelength region of an electromagnetic spectrum beyond visible range does not have sufficient quantum energy to ionize an atom, therefore it is known as non-ionizing radiation. However, while interacting with matter, this radiation can produce thermal effects and change the rotational and vibrational configurations of molecules and atoms. This chapter deals with the spectroscopic methods based on such vibrational and rotational transitions occurring in longer wavelength region such as infrared (IR) spectroscopy, nuclear magnetic resonance (NMR), and electron spin resonance (ESR). It also includes mass spectroscopy (MS) based upon segregation of ions according to their mass-to-charge ratio.
9.2
INFRARED SPECTROSCOPY Infrared spectroscopy is based on the excitation of molecules on absorption of electromagnetic radiations in the IR region. This region lies between visible and microwave region of the electromagnetic spectrum, and is characterized by longer wavelength ranging from 1 µm to 300–400 µm. The IR region is divided into three regions: the near-, mid-, and far-IR region (Figure 9.1). The region closest to the visible region, called the near-IR region, possesses energy that can bring about harmonic or anharmonic molecular vibrations. The mid-IR region can cause some fundamental vibrations and rotations, and the far-IR region having low energy can only cause rotational transitions. All these vibrations and rotations occur at a specific frequency characteristically when absorbed by the molecules according to their molecular
260 Bioanalytical Techniques
Figure 9.1
Infrared region
structures. These frequencies are in accordance with the transition energy of the bond or group that gives rise to an IR spectrum which is the signature of that molecule, known as the fingerprint of the IR region.
9.2.1
Basic Principles
A molecule consisting of n atoms is said to have 3n degrees of freedom. A non-linear molecule possesses 3 degrees for rotation and 3 for translation, therefore, it is left with (3n – 6) degrees for vibrational modes. For example, a non-linear molecule like CH4 has 9 degrees of freedom for vibrations. A linear molecule like CO2 has 2 degrees for rotation and 3 for translation, therefore, it has (3n – 5) degrees for fundamental vibrations and can vibrate only in 4 modes. Two types of molecular vibrations, stretching and bending, are illustrated in Figure 9.2.
Figure 9.2
Types of bond vibrations
Spectroscopy II
261
The basic requirement for absorption of IR radiation is that there must be a net change in dipole moment during the vibration for the molecule or the functional group. A highly symmetrical molecule such as O2, N2, or ethylene (CH2CH2), in which the stretching of a bond does not produce any change in the dipole moment, cannot be excited by IR radiation, and, therefore, does not show any IR peak in the spectrum. Similarly, the stretching vibrations of completely symmetrical double and triple bonds, like CH3C∫CCH3, does not result in the change in the dipole moment, and, therefore, are IR inactive (Figure 9.3). Polar groups like water absorb IR light more strongly and give an intense peak because of high dipole moment. The greater the polarity of the bond, the stronger is the IR absorption. For example, the carbonyl bond is very polar, and absorbs very strongly while the C∫C bond in most asymmetric alkynes is much less polar, therefore, absorbs rather weakly. A molecule can only stretch and bend at certain defined frequencies and if it is exposed to electromagnetic radiation of this frequency, it will absorb energy from the radiation and jump to a higher vibrational energy state. For most of the organic compounds, these frequencies correspond to the IR region. The transitions between vibrational energy levels caused by IR absorption are said to be equivalent to a classical harmonic oscillator and are expressed by Hooke’s law:
H3C O
C
O
CH3 C
C
H3C
H3C
C
C
CH3
CH3 Infrared-inactive double and triple bonds (a)
H
H
H
H
O
O
Symmetric
Asymmetric
H
H O Bending
(b)
Figure 9.3
Infrared-inactive (a) double and triple bonds and (b) water molecules showing various stretching
262 Bioanalytical Techniques __
1 k n = ____ __ 2pc m where n is the fundamental vibration frequency, k is the force constant, and m is the reduced mass. m can be given by the following relation:
÷
m1m2 m = _______ m1 + m2
9.1
9.2
where m1 and m2 are the component masses for the chemical bond under consideration. This relationship provides a reasonably good fit between the bond stretching vibrations predicted in a molecule and the values observed for the fundamental vibrations. Any anharmonicity causing deviations from this fit accounts for weaker absorptions, known as overtones, resulting from the transition occurring from the ground vibrational state (n = 0) to the second, third, fourth vibrationally excited state (that is, n = 2, 3, or 4) rather than the first (n = 1). The frequencies of first (n = 0 Æ n = 2) and second (n = 0 Æ n = 3) overtone bands are 2–3 times the frequency of the fundamental vibrations and appear in the region of shorter wavelengths (Figure 9.4). Band intensity in an IR spectrum is expressed as absorbance (A) or transmittance (T), and the relation between them can be given as A = log10 (1/T)
9.2.2
9.3
Infrared Spectrophotometer
An IR spectrophotometer uses a radiation source for emitting radiations in the IR region. The radiation sources commonly used are a Nernst glower made of n=3 n=2 n=1 n=0 Fundamental
Figure 9.4
First overtone
Second overtone
Energy-level diagram depicting transitions for fundamental and overtone IR bands
Spectroscopy II
263
zirconium oxide heated to about 1500°C–1700°C to give radiation at 7100 cm –1 or a globar source made of silicon carbide heated to about 1500°C for emission at 5200 cm–1. To obtain radiation in far-IR region (>200 cm–1), mercury arc lamps are usually employed. A monochromator is used for splitting the source radiation into different wavelengths, followed by a slit to allow the radiation to pass through the sample. A double-beam IR spectrophotometer consists of a beam splitter that separates the incident radiation into two beams: one moving towards the sample, and other to the reference (Figure 9.5). Sample cells are made of IR transparent material. Solid samples are generally prepared by mulling in Nujol or by grinding with potassium bromide (KBr), and then pressing into discs. Although most mulling agents show absorption bands in the spectrum which may coincide with the sample bands at same frequency in the spectrum. Liquid samples are easy to treat and are simply made into thin films of just a few microns of thickness, between the plates of NaCl, KBr, or another suitably transparent material until absorption bands of measurable intensity are obtained. Water and alcohols are not used as solvents since they absorb strongly in the IR region due to O–H vibration. Gaseous samples are filled in multi-reflection gas cell having long path length. In this cell, a series of mirrors are placed which deflect the IR beam back and forth many a times until it exits in the cell after having travelled the required path length. The radiations that pass through the sample are collected by the detector in comparison to those passing through the reference. The output is generated as a function of frequency or wavelength.
Figure 9.5
Diagrammatic representation spectrophotometer
of
an
instrumental
set-up
for
IR
264 Bioanalytical Techniques FOURIER TRANSFORM INFRARED SPECTROSCOPY Infrared spectroscopy has been dramatically improvised by the development of the Fourier transform method after the invention of an interferometer by A. A. Michelson in the 1890s. However, it was not until the late 1960s when microcomputers able to do the Fourier transformation became available and commercial Fourier transform infrared (FTIR) spectrometers appeared. It is a non-destructive technique which provides precise measurement with rapidity and high sensitivity. Due to this reason, FTIR spectroscopy is now preferred over dispersive IR methods of IR spectral analysis. See Figure A.
Figure A
Schematic representation of FTIR spectrometer
An FTIR spectrophotometer is an interferometer. The radiation from an IR source is directed to the sample cell. The sample absorbs all the wavelengths characteristic of its spectrum. A beam splitter splits the light into two component beams, one of which is reflected from a fixed mirror while the other is reflected from a mirror which moves continuously over a distance of about 2.5 micrometres (mm). As the two beams recombine at the detector, an interference pattern is produced due to difference in path lengths which generates an interferogram. A single scan of the entire distance takes about 2 seconds and multiple scans from n points are stored in the computer interfaced to the instrument. This process improves the signal-to-noise ratio significantly. The data obtained in the form of intensity versus time spectrum is converted by the computer into intensity versus frequency spectrum using the Fourier transform mathematical function. The linear relationship between the absorbance and the number of absorbing molecules may also enable multicomponent analysis of compounds.
Spectroscopy II
9.2.3
265
Infrared Spectrum
Absorbance
Transmittance
Infrared spectroscopy reveals a lot of information for the identification of organic as well as inorganic compounds. It can be used as a fingerprint for structural determination of backbone of the molecule and its functional groups on the basis of the absorption of particular IR frequencies corresponding to the vibration of specific sets of chemical bonds within a molecule. The absorption depends upon several factors such as nature of bonds, polarity of functional groups, unsaturation (double or triple bonds), cis- or trans-positions of functional groups, and aromatic structure. The information deduced by the absence of some characteristic bands is also vital for ruling out the presence of certain groups or bonds, helping greatly in presuming the molecular structure and composition. An IR spectrum is typically a plot of transmittance versus wavenumber (cm–1). The solid line traces the values of per cent transmittance for every wavelength so that the regions of strong absorption appear as downward pointing peaks (Figure 9.6). Infrared bands at frequencies 4000–1250 cm–1 are generally associated with changes in vibrational energies and are characteristic types of bonds present in a molecule. Absorption of lower IR frequencies (less than 1250 cm –1) is usually related to complex vibrational and rotational energy changes in a molecule which are highly characteristic for such changes. This region is, therefore, called the “fingerprint” region, where the IR spectrum of even the most closely related molecules is different. The fundamental vibrations due to C–H, O–H, and N–H stretching appear in the 4000–2500 cm−1 region. Aliphatic C–H stretching bands are found
3800 3400 3000 2600 2200 1800 1400 1000 –1 Wavenumber (cm ) (a)
Figure 9.6
3800 3400 3000 2600 2200 1800 1400 1000 –1 Wavenumber (cm ) (b)
An IR spectrum plotted for (a) absorbance and (b) transmittance
266 Bioanalytical Techniques at 3000–2850 cm−1, C–H bond in an aromatic ring or a C–H bond near a double bond shows peak between 3100 and 3000 cm−1. A broad band at 3700–3600 cm−1 is a characteristic of O–H stretching, while a sharp band at 3400 cm–1 and 3300 cm−1 is a characteristic of N–H stretching. IR bands in the region of 2500–2000 cm−1 characterize triple-bond stretching, of which the C∫C bonds absorb between 2300 and 2050 cm−1, while the C∫N bonds absorb between 2300 and 2200 cm−1. Double-bond stretching is seen in the region of 2000–1500 cm−1, of which C=O bond absorbs at 1830–1650 cm−1, while weak C=C stretching and strong C=N stretching occur approximately at 1650 cm−1 (Figure 9.7). Infrared spectrum obtained for an “unknown” compound can be compared with previously recorded reference spectrum.
9.2.4 Applications of Infrared Spectroscopy The most common use of IR spectroscopy is identification of organic and inorganic IR active compounds. The IR spectra of inorganic molecules are determined by the normal modes of vibration exhibited by such molecules. IR spectroscopy is widely used for characterizing compounds in coordination chemistry, particularly coordination modes of a ligand and the geometrical arrangement of ligands around the metal atom. It is a popular method for identifying structures of polymers as well as their cross-linking properties. This technique is frequently used for investigating the structural properties of clay minerals. Biological systems including lipids, proteins, peptides, biomembranes, nucleic acids, animal tissues, microbial cells, plants, and clinical samples have all been successfully studied by using IR spectroscopy.
Transmittance (%)
100
0 4000
3000
2000
1500
1000
1500
Wavenumber (cm–1)
Figure 9.7
Infrared spectrum showing characteristic peaks of some functional groups
Spectroscopy II
9.3
267
NUCLEAR MAGNETIC RESONANCE SPECTROSCOPY Nuclear magnetic resonance is an extremely versatile technique used for structural elucidation of wide range of chemical and biochemical compounds. It can be applied for functional group analysis, kinetic studies of enzymatic and non-enzymatic reactions, determination of the stereochemistry of compounds, and study of molecular conformations of complex macromolecules such as DNA and proteins. This technique relies on the absorption of a characteristic frequency of electromagnetic radiation of radio wave region (3 kHz–300 GHz) by nuclei having odd number of protons or neutrons on the application of or exposure to magnetic field. This gives the technique its name: proton NMR or 1H NMR. The resonant frequency of nuclei is proportional to the magnetic field strength and depends upon the nuclear magnetic moment.
9.3.1
Basic Principles
Nuclei of atoms possess the spinning characteristic that generates small magnetic field which results in a magnetic moment proportional to the spin. Each spin exists in multiples of ½ and + ½ or – ½ spin can be assigned to each nucleus. If spinning nuclei having opposite spins pair up, their spins cancel each other so that they have no overall spin. The nuclei having unpaired particles (such as 1H and 13C) possess an overall spin and have an instrinsic nuclear angular momentum. NMR is concerned with these unpaired nuclear spins; therefore, NMR can only be performed on naturally abundant isotopes having nuclear spins. The spin angular momentum possible for a spinning nucleus can be given as Spin angular momentum = [I (I + 1)]1/2 h/2p
9.4
where I is the spin quantum number of the nucleus and h is the Planck’s constant. The magnetic moment µ of the nucleus is given as m = g (spin angular momentum). Therefore, m = g [I (I + 1)]1/2 h/2p 9.5 where g is the proportionality constant called gyromagnetic ratio. The overall spin of a nucleus can be determined by any of the three factors given here:
268 Bioanalytical Techniques (i) If the nucleus has even number of protons and even number of neutrons, the net spin will be zero. (ii) If the nucleus has odd number of neutrons and odd number of protons, it will have an integer spin, that is 1, 2, and 3. (iii) If the sum of the number of neutrons and the number of protons in the nucleus is an odd number, it will have a half-integer spin that is 1/2, 3/2, 5/2, etc. These rules can be summarized in terms of atomic mass and atomic number as shown in Table 9.1. When a nucleus having some magnetic moment is placed in an external magnetic field Ho, its spin vector aligns itself with the external field. If the vector alignment is parallel to the direction of the external field, the configuration is said to be in the low-energy state (spin state = +1/2), but if the vector alignment lies in the anti-parallel direction to that of the external field, the configuration it is said to be in the high-energy state (spin state = –1/2). A spinning nucleus rotates (spins) about its own axis and precess about the axis of the magnetic field H. On absorption of photon energy, the nucleus “flips” its spin so that it opposes the applied field and enters the higher energy state. This transition is possible if the energy of this photon is at least equal to the energy difference between the two energy levels, which depends upon the strength of the external magnetic field used. This energy requirement for nuclear transition is very low and is conveniently sufficed by radio waves (Figure 9.8). The difference in energy between the transition levels can be calculated as DE = E2 – E1 = (g h H) / 2p
9.6
Table 9.1 Spin quantum number (I) of NMR active and inactive elements Examples 12 1
H,
35
19
Cl,
13
C
17
O
2
16
C,
H,
O,
F,
79
14
32
31
P
Br,
N
S
127
I
Number of protons Number of neutrons
Spin quantum number (I)
Response to NMR
Even
Even
0
NMR inactive
Odd
Even
1/2
NMR active
“
“
3/2
NMR active
Even
Odd
1/2
NMR active
“
“
5/2
NMR active
Odd
Odd
1
NMR active
Spectroscopy II
Energy
No field
Applied magnetic field
0 Precessional orbit
Figure 9.8
269
Spinning nucleus
Energy-level diagram showing spin states of a nucleus under applied magnetic field
where E2 is the energy of nucleus after transition to high-energy state, E1 is the energy of nucleus before transition from low-energy state, g is the gyromagnetic ratio, h is the Planck’s constant, and H is the applied magnetic field. This relationship suggests that if the magnetic field, H, is increased, DE also increases. Also if a nucleus has a relatively large gyromagnetic ratio, then DE will also be proportionately larger. On absorption of a perfect radio frequency (RF) that causes flipping of the proton from one magnetic alignment to another, a resonance is established. This frequency “n” is called the Larmor frequency and is given by g 9.7 n = ___ H 2p where n is the Larmor frequency, g is the gyromagnetic ratio, and H is the applied magnetic field. Since the nuclei are not the isolated entities and remain surrounded by number of electrons depending upon the atomic structure, each nucleus experiences different effect of external magnetic field. At a particular RF, each nucleus or proton needs a slightly different magnetic field to be applied to it in order to bring it into the resonance condition. This resonance is detected and the signals are converted into characteristic peaks in an NMR spectrum.
270 Bioanalytical Techniques However, generally in a sample only a few nuclei are in the lower energy state while others exist in higher energy state owing to their respective spin orientations. Group of spins experiencing the same magnetic field strength is referred to as a spin packet. NMR relies on the net magnetization from all of the spin packets which is the vector sum of all the individual magnetization vectors. On providing the energy, all nuclei reorient in the magnetic field which can bring about a change in the net magnetization. If the number of nuclei in the higher energy levels becomes equal to those in the lower energy levels, no further absorption of radiation takes place and the spin system is said to be saturated where net magnetization becomes zero. The excited nuclei lose their energy of excitation and return to the unexcited state; the process is called relaxation and the time spent in the excited state is the relaxation time.
9.3.2
Instrumentation
The instrumental set-up of NMR is extensive, consisting of a complex network of components. The primary necessity is to produce the magnetic field (Ho) to bring about the nuclear resonance in the sample. This field is produced through a superconducting magnet. Usually, the field strength varies from 6 to 23.5 Tesla approximately. The resonance frequency is kept constant by using a field frequency lock, tuned to a reference NMR resonance frequency called lock frequency. Usually NMR resonance frequency of deuterium is used as a lock frequency, as it is used as a solvent to prepare the sample. Since the field lock is based on the deuterium resonance which is very temperature sensitive (giving a resonance shift of 0.01 ppm/K), a very stable temperature control in the probe is maintained to ensure the optimal performance of the NMR spectrometer. This is done using a temperature control system which maintains the temperature within ±0.05 K of the set temperature. The magnet is accompanied by shim coils which create a small magnetic field for homogenizing the applied magnetic field. These shim coils consist of a sample probe which holds the sample in a spinner. A pulse of RF radiation is generated using RF coils to rotate the spins and resonate the nuclei at Larmor frequency. The RF coils also receive and detect the signals from the sample. All the components of the NMR spectrometer are controlled by the computer. The RF source is programmed to produce a sine wave of the desired frequency. The pulse programmer is used to set the width and shape of the RF pulses. The weak NMR signals are amplified and the spectrum is generated on the display (Figure 9.9).
Spectroscopy II
Figure 9.9
9.3.3
271
Schematic representation of NMR spectrometer
Interpretation of Nuclear Magnetic Resonance Spectrum
The 1H NMR spectrum of a compound provides information regarding different types of protons present in the molecule, which is sufficient to elucidate the structure of the compound. (i) Number of signals The number of signals in an NMR spectrum gives information about the number of different kinds of protons experiencing different chemical environments present in the molecule. A set of protons which are in the same magnetic environment or, in other words, experiencing same magnetic force are termed as equivalent protons. The number of signals on the NMR spectrum thus represent the number of different sets of equivalent protons, since every chemically distinct proton will give a distinct peak pertaining to its resonance in the NMR spectrum. Thus, different signals in an NMR spectrum are indicative of the presence of non-equivalent groups of protons. See Figure 9.10. (ii) Integration of signals In 1H NMR spectroscopy, the area under an NMR signal is proportional to the number of protons resonating at the frequency to which the peak corresponds. The areas represented by the integration step function of the computer system of NMR instrument are usually integrated and displayed as numerical values under the d scale, directly indicating the relative number of chemically distinct protons. However, it
272 Bioanalytical Techniques
2,2-dichloro propane has only one type of protons and gives only one peak
Cl H3C
C
CH3
Cl 10
9
8
7
6
5 4 ppm
3
2
1
0
8
7
6
5 4 ppm
3
2
1
0
8
7
6
5 4 ppm
3
2
1
0
HPM-00-189
2-chloropropane has two types of protons and gives two peaks
(A) H H3(C)
C
CH3
Cl 10
9
HPM-03-171
1-chloropropane has three different types of protons and therefore gives three peaks
CH3
CH2
CH2
Cl
10
9
HPM-03-17
Figure 9.10
NMR signals for (a) 2,2-dichloro propane (b) 2-chloropropane
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273
does not give the absolute number of protons. The integration function is useful in determining the relative amounts of two or more compounds in a mixed sample. (iii) Chemical shift The resonance frequency for a proton is determined by the strength of the external magnetic field. However, the effective magnetic field (H) at the nucleus is not equal to the applied magnetic field (Ho) and is affected by the molecular environment of the nucleus because of the shielding effect of surrounding electrons. This is observed as shifts in positions of NMR signals and is referred to as “chemical shifts”. Chemically equivalent protons have the same chemical shift and, therefore, give rise to the same signal. The relationship is given as H = Ho (1 – s)
9.8
where s is the shielding constant. Due to shielding, the nucleus requires greater applied field strength for resonance, called diamagnetic shift or upfield shift. If the effective magnetic field is larger, the nucleus is said to be deshielded causing a paramagnetic shift or a low field shift. The chemical shift (d) is very useful in determining the chemical environment around a nucleus and can be measured by comparing the resonance frequency of the sample to that of a reference compound, and is given as d = (uSample – uRef ) / uRef × 106
9.9
where d is the chemical shift, uSample is the resonance frequency of sample, and uRef is the resonance frequency of reference compound. The chemical shift is dimensionless and is reported in ppm. The value of chemical shift not only depends on the molecular structure, but also on the experimental conditions such as type of solvent, temperature, and strength of magnetic field. For ease of comparison of 1H NMR spectra, a chemical standard named tetramethylsilane or TMS Si(CH3)4 is commonly used. It is an inert compound which does not react with the sample. It has 12 equivalent protons which absorb at the same magnetic field and result in intense absorption. The chemical shift of TMS is taken as 0.00 ppm and most organic compounds have chemical shifts between 0 and 15 ppm.
274 Bioanalytical Techniques Table 9.2 Chemical shifts of some classes of organic compounds Class of compound
General formula
Chemical shift (d) ppm
Primary alkane
R-CH3
0.9
Secondary alkane
R 2CH2
1.3
Tertiary alkane
R3CH
1.5
Alkyl chloride
>CHCl
3–4
Alcohol
>CHOH
3.4–4
Ether
>CHOR
3.3–4
Aromatic
Ar-H
6–8.5
Phenolic
Ar-O-H
4–12
(iv) Splitting of signals Often NMR spectra are complicated and there are cases where simply the number of non-equivalent protons does not comply with the number of signals. On further resolving the low-resolution NMR spectrum, the signals may split into several lines. Non-equivalent protons which are close to one another exert an influence on each other’s effective magnetic field which is observed as splitting in signals if the distance between them is less than or equal to three sigma bonds. This is called spin–spin coupling or J-coupling. The magnitude of the observed spin splitting is independent of external magnetic field and is given by the coupling constant J (in Hz) (Figure 9.11). The complexity of the splitting pattern in a spectrum increases as the number of neighbouring non-equivalent protons increases. Therefore, the amount of splitting directly indicates the number of hydrogen atoms attached to the carbon atom or atoms in vicinity to the one under consideration. If there are no protons in the vicinity, a single peak will appear, known as a singlet. If there is one proton in the vicinity, the resonance will be split into two peaks of equal size appearing as a doublet. If there are two protons in the vicinity, three peaks will appear as a triplet with their area in the ratio of 1:2:1. Similarly, if there are three protons in the vicinity, a quartet of four peaks will appear with their area in the ratio of 1:3:3:1. In general, an NMR resonance will split into n + 1 peaks (n = number of protons on the adjacent atom), following the Pascal’s triangle. For example, ethyl acetate spectrum displays the typical quartet and triplet of a substituted ethyl group. In the NMR spectrum of ethanol, the three different types of H atoms give three signals. The hydrogen of OH group is not coupled with the other H atom and, therefore, occurs as a singlet, but the hydrogen of CH3 and >CH2
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275
C No coupled hydrogen
C
C
C
H A singlet
C J H One coupled hydrogen
C
C
C
H A doublet
C H Two coupled hydrogens
H
C
J C
J
H A triplet
C H Three coupled hydrogens
H
C
J C
J
H
H
J A quartet
(a)
1 1 1 1 2 1 1 3 3 1 1 4 6 4 1 1 5 10 10 5 1
Coupling to first nucleus (N=1, doublet) Coupling to second nucleus (N=2, triplet) Coupling to third nucleus (N=3, quartet) Coupling to fourth nucleus (N=4, quintet) Coupling to fifth nucleus (N=5, sextet) (b)
Figure 9.11
(a) Splitting of NMR signals due to J-coupling (b) Pascal’s triangle showing relative peak intensities for multiplet peaks arising from spin– spin coupling of a nucleus to N equivalent nuclei
are coupling with each other, splitting the signals into a triplet and quartet, respectively.
9.3.4 Applications The NMR spectrum provides valuable information about the molecular structure of a compound. The versatility of the technique has made it a wonderful tool in organic and inorganic chemistry for widespread applications. It is now routinely used for the determination of the structure of complex molecules like proteins. This technique allows the characterization of biomolecules at atomic resolution.
276 Bioanalytical Techniques It helps to study protein folding and unfolding patterns, and in determining the intermediate and residual structures of unfolded proteins. NMR is also helpful in studying intermolecular interactions such as protein–ligand binding, studying conformations of the compounds bound to enzymes or receptors, studying active sites and domains of proteins, etc. (see Figures 9.12 and 9.13).
9.4
ELECTRON SPIN RESONANCE SPECTROSCOPY Electron spin resonance spectroscopy, also called electron paramagnetic resonance (EPR) spectroscopy, is meant for studying chemical species with
Figure 9.12
NMR spectrum of ethyl acetate
H
H
H
C
C
H
H
OH
H
H
H
4
Figure 9.13
3
NMR spectrum of ethanol
2
1
0
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277
unpaired electrons such as stable inorganic paramagnetic molecules, organic molecules carrying odd number of electrons, free radicals, transition metal complexes, rare earth elements, and others. Electron spin resonance spectroscopy is based on the resonance of unpaired electron by absorption of microwave radiation, measuring the transition frequency between the two electron spin states. This technique is analogous to NMR in many aspects as in NMR similar resonance phenomenon occurs in nucleus having odd number of protons or neutrons on absorption of radio waves. However, it differs from NMR since the resonant energy in ESR corresponds to the frequency in the order of GHz as compared to the MHz frequency encountered in NMR because the energy difference between the two spin states of an electron in a reasonably strong magnetic field is quite large.
9.4.1
Basic Principles
The two electron spin energy states remain degenerated in the absence of external magnetic field. However, under an applied magnetic field (Ho), the energy level splits into two: a low-energy state corresponding to the spin parallel to the magnetic field vector where the electron has spin angular momentum ms equal to –1/2, and a high-energy state corresponding to the spin anti-parallel to the magnetic field vector where electron has spin angular momentum ms equal to +1/2. The split in energy levels due to the magnetic field is called the Zeeman interaction, which is observed as splitting of spectral lines (Figure 9.14). The transition of electron between the two energy levels takes place by the absorption of radiation of suitable frequency, known as the resonant frequency, which lies in the microwave region (9–10 GHz).
Energy
ms = 1/2
DE = hu = gbH
ms = –1/2 Bo = O
Figure 9.14
Applied magnetic field
Energy levels for an electron spin (mS = ±1/2) in an applied magnetic field H
278 Bioanalytical Techniques The transition energy (DE) is given by the following relation: DE = hu = gbH
9.10
where h is the universal Planck’s constant, u is the frequency of microwave radiation, b is the Bohr magneton, H is the strength of magnetic field, and g is the splitting factor. In the above equation, the two variables are u and H which indicate that the transition energy is directly proportional to the strength of applied magnetic field at the selected frequency. The splitting factor (g) is a proportionality constant, also called Lande’s splitting factor, dependent upon the surrounding environment of the electron. Its value for a free electron is 2.0023, which is same for the electrons in free radicals. In gases the molecules are in free motion, therefore, their g-value is an average of all their orientations but it differs in ionic crystals. It affects the ESR signal, therefore, the effective g-value is calculated by using a reference of known g-value. A very common reference is the stable free radical diphenylpicrylhydrazyl (for which g = 2.0036). The net absorption of energy from microwave radiation depends upon the equilibrium distribution of the electrons between the two spin energy states, which is given by Maxwell–Boltzmann distribution: hn DE hhigh / hlow = exp –____ = exp – ___ KT kT
(
)
(
)
9.11
where hhigh is the number of electrons in higher energy state, hlow is the number of electrons in lower energy state, k is Boltzmann constant, and T is the temperature in kelvins. According to Maxwell–Boltzmann distribution, the population of electrons in low-energy state is usually higher than in high-energy state. If at a constant temperature, the strength of applied magnetic field is increased, the difference between the two populations can also be increased. On supplying microwave radiation of a suitable frequency, the electrons can be excited from the low-energy electrons to the high-energy state.
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279
Greater sensitivity in ESR can be achieved by keeping the temperature low, and working at high frequencies and proportionately high magnetic field strength. On applying the field of about 3000 G, at T = 298 K, X-band microwave frequencies (n ª 9.75 GHz) give hhigh /hlow = 0.9986. Although this value is near to 1, meaning that the two-spin energy levels of electrons are almost equally populated, still there occurs net absorption of energy which brings about the transitions from the lower to the higher level. However, if the radiation of slightly high intensity than required is incident, it tends to equalize the populations causing a decrease in net absorption, which is referred to as saturation. Further, in order to maintain a population excess in the lower energy level, the electrons from the upper energy level give up the energy to their surroundings to return to the lower level. The process is called spin relaxation process, which may occur by spin–lattice relaxation or spin–spin relaxation. In spin–lattice relaxation, the energy is dissipated to the neighbouring molecules within the lattice, whereas in spin–spin relaxation the exchange of energy occurs between high-energy electrons and the electrons in vicinity, without any energy transfer to the lattice. The ESR spectrum is obtained by plotting the intensity or absorbance against the magnetic field strength. Usually, it is represented as a first derivative of the absorption curve (Figure 9.15).
9.4.1.1
Hyperfine interactions
Signal
Signal
The magnetic moment of spinning electron is also affected by the magnetic moments of nuclei spinning in vicinity. Such interactions are called hyperfine interactions and cannot be taken for granted as they cause further splitting of the two spin energy levels of the electron giving rise to additional lines in
First derivative
Absorbance Magnetic field strength
Figure 9.15
An absorption curve and its first derivative
Magnetic field strength
280 Bioanalytical Techniques
Figure 9.16
Hyperfine splitting of the energy states of an electron
the ESR spectrum (Figure 9.16). In fact these interactions provide additional information about the micro-environment of unpaired electrons. The number of lines (Nhyperfine) from the hyperfine interaction can be determined by the formula given below and the intensity of each line follows Pascal’s triangle. Nhyperfine = (2nI + 1)
9.12
where n is the number of equivalent nuclei and I is the spin quantum number. For example, methyl radical (CH3) has three 1H nuclei each with I = 1/2, and so the number of lines expected is 2nI + 1 = 2(3)(1/2) + 1 = 4. Thus, the signal splits into four lines, the ratio of each line’s intensity is 1:3:3:1, corresponding to the Pascal’s triangle (Figure 9.17).
9.4.2
Electron Spin Resonance Instrumentation
The detection of ESR absorption in a system containing unpaired electrons requires a spectrometer with a high and extremely stable magnetic field. An ESR spectrometer operates at a fixed microwave frequency and scans ESR spectrum by a linear variation of the static magnetic field. A klystron is the source of microwave consisting of a heated cathode that emits electrons, an anode to collect the electrons, and a reflector electrode at high negative potential which reflects the electrons back to the anode. The motion of the electrons
281
Derivative signal
Spectroscopy II
Field strength
Figure 9.17
EPR spectrum of methyl radical
creates an oscillating electric field which produces the desired electromagnetic radiation. The microwave frequency can be tuned by adjusting the distance between the anode and the reflector cathode or by adjusting the reflector voltage. Klystron is accompanied by an isolator and an attenuator. The isolator prevents scattering of radiation as well as eliminates the vibrations, and an attenuator, made of resistive material, is used to adjust the microwave power. The microwaves from the Klystron are routed through a circulator towards the resonant cavity containing the sample. The length of the cavity is exactly one wavelength (X-band cavity is 3 cm long corresponding to its wavelength), so that the Klystron frequency is tuned to the “cavity resonant frequency”. The magnetic field is generated by a large electromagnet and can be varied from 0 to 500 Gauss. The homogeneity and stability of the field is maintained through controlled power supply which also allows the magnetic field to be swept over a width to achieve the resonance condition. Silicon crystal detectors are commonly used for converting the microwave radiation into a direct current (DC) output. The modulation of the signal at a frequency complying with a good signal-to-noise (S/N) ratio in the crystal detector is accomplished by introducing an alternating current (AC) signal to the modulation coil. In order to adjust the spectrometer and to observe the signal, a cathode ray oscilloscope is employed. An X–Y recorder is used for recording the signal of an ESR spectrum which is usually displayed as a first derivative of the absorption line to improve the S/N ratio.
282 Bioanalytical Techniques
Figure 9.18
Schematic representation of an ESR spectrometer
9.4.3 Applications of Electron Spin Resonance Electron spin resonance is a powerful, non-destructive, and non-intrusive analytical method providing information about the paramagnetic species in the sample. The spectral lines suggest the type of nuclei present in the vicinity of the unpaired electron. ESR yields useful structural information which can be used in a variety of ways to deduce the probable molecular structure and composition of the sample. This technique also allows the study of enzymatic kinetics. One of the most impressive applications of ESR is the determination of free radicals present in extremely low concentrations (Figure 9.18).
9.5
MASS SPECTROSCOPY Mass spectrometry is an analytical technique used for identification and quantification of wide range of chemical and biochemical compounds. The sensitivity of this technique allows detection of minor mass changes in the compounds, allowing for their structural elucidation. In MS, molecules are converted into gaseous phase ions and separated on the basis of their massto-charge ratio (m/z). It is based upon the principle that motion of charged particles (ions) under an electric or magnetic field is affected by the m/z of the
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283
ions and the result is obtained in the form of a mass spectrum which represents the relative abundance versus m/z of a sample. The ionization methods used commonly for majority of biochemical analyses are electrospray ionization (ESI) and matrix-assisted laser desorption ionization (MALDI). With most ionization methods there is the possibility of creating both positively and negatively charged sample ions, depending on the proton affinity of the sample.
9.5.1
Basic Principles
The analyte molecules are converted to gaseous phase and are bombarded with high speed electrons for their ionization. M + e – Æ M+ + 2e –
9.13
These ions are separated in space or time, based on their m/z and are quantified. All the ions with equal charge, accelerated across the same distance by equal amount of force, possess equal kinetic energy (KE), determined by the acceleration voltage of the instrument (V ) and the charge of the ion (z), given as K1E1 = zV = ½ mv2
9.14
where, m is the mass of ion and v is the velocity of ion. As the ions move through the magnetic field, their velocity is given as the ratio of distance travelled by the ion and the time taken, that is v = D/t
9.15
where v is the velocity of the ion, D is the distance travelled by the ion, and T is the time taken by the ion to travel the distance. On substituting v in Equation 9.14, we get zV = ½ m (D/t)2
9.16
On rearranging Equation 9.16, we get ___
____
t = d/÷2V (÷m/z )
9.17
284 Bioanalytical Techniques D is fixed by instrument design and V can be held constant electrically, that is ____
t = k ÷m/z
9.18
Therefore, m/z is proportional to the square of the travel time t.
9.5.2
Mass Spectrometer
The main components of a mass spectrometer are as follows (Figure 9.19): (i) Sample inlet The gaseous phase sample ions are introduced into the ionization source of the instrument through sample inlet. Common method for sample introduction is direct vapour inlet, in which the high vapour pressure sample in gaseous phase is introduced directly into the source region through a needle valve. If the analyte is thermally labile, that is, it decomposes at higher temperatures, or if it does not have sufficient vapour pressure, the sample must be directly ionized from the condensed phase. The sample can also be introduced to the ionization source through a chromatographic system such as a coupled high pressure liquid chromatography or gas chromatography instrument. (ii) Ionization source The sample molecules are excited so that they eject an electron to form a radical cation (M+) or are forced to undergo ion– molecule reactions to produce adduct ions (MH+), mostly by addition of a proton (H+) to the molecule (M). M + H+ = MH+
9.19
The ionization energy is however significant because it controls the amount of fragmentation observed in the mass spectrum. Although this fragmentation complicates the mass spectrum, it provides structural information for the identification of unknown compounds. Some ionization techniques are very soft and only produce molecular ions. Other techniques
Figure 9.19
Schematic representation of a mass spectrometer
Spectroscopy II
Figure 9.20
285
Electron ionization technique
are very energetic and cause ions to undergo extensive fragmentation, known as hard ionization. The ionization method to be used depends upon the type of sample under investigation and the mass spectrometer available. Electron ionization and chemical ionization are only suitable for gas phase ionization. Fast atom bombardment (FAB), secondary ion mass spectrometry, electrospray, and MALDI are used to ionize condensed phase samples (Figure 9.20). (iii) Analysers As soon as the ions are formed in the source region, they are accelerated into the mass analyser by an electric field. The main function of the mass analyser is to resolve the ions formed after ionization according to m/z. These mass analysers have different features, including the range of m/z analysed, the mass accuracy, detection limits, resolution, and scan rates (Figure 9.21). Analysers can be categorized as continuous or pulsed. Continuous analysers transmit a single selected m/z to the detector and the mass spectrum is obtained by scanning the analyser so that ions with different m/z are detected. These analysers include quadrupole filters and magnetic sectors. These work on the principle of single ion monitoring (SIM) which improves the S/N ratio. Pulsed mass analysers collect an entire mass spectrum from a single pulse of ions. Since all the ions are detected, these analysers have very high transmission efficiency which increases the S/N ratio. Pulsed analysers include time-of-flight (TOF), ion cyclotron resonance, and quadrupole ion trap mass spectrometers. The compatibility of different analysers with different ionization methods varies. For example, all of the analysers listed above can be used in
286 Bioanalytical Techniques
Figure 9.21
Schematic representation of (a) magnetic sector mass analyser, (b) quadrupole mass analyser, and (c) quadrupole ion trap mass analyser
Spectroscopy II
287
conjunction with electrospray ionization, whereas MALDI is usually coupled to a TOF analyser. Tandem mass spectrometers consisting of two mass analysers have been developed for sophisticated structural and sequencing studies of complex macromolecules like proteins. In these instruments, the first analyser is used as a mass-filter to isolate a molecular ion that will enter the collision cell where it is fragmented. The second mass analyser analyses the resulting fragments measuring their molecular mass. Commonly coupled tandem mass spectrometers include quadrupole–quadrupole, magnetic sector– quadrupole, and the quadrupole–TOF analysers. (iv) Ion detectors The ions sorted by the analyser strike the detector where the signals are amplified and transmitted to the data system. The m/z values of the ions are plotted against their intensities to generate the mass spectrum which shows the number of components in the sample, the molecular mass of each component, and the relative abundance out of the various components in the sample. Common detectors used in mass spectrometers are the photomultiplier tube, the electron multiplier tube, the micro-channel plate detectors, and Faraday cup. (v) High vacuum system Since mass spectrometers use gas-phase ions for analysis, they operate in a high vacuum system typically 10 –2 to 10 –5 Pa, often maintained by a pumping system. The ionization source, analyser, and detector of the mass spectrometer are maintained under high vacuum so that the ions travel smoothly through the analyser to the detector without any hindrance. This also minimizes the chance of collision of ions with other molecules so that they do not react, scatter, or break in fragment.
MATRIX-ASSISTED LASER DESORPTION IONIZATION-TIME-OF-FLIGHT MASS SPECTROMETRY Matrix-assisted laser desorption ionization is a soft ionization technique used nowadays for analysis of thermolabile, non-volatile organic compounds of high molecular weight (about 50 kDa) such as proteins, peptides, and DNA fragments. The sample is co-crystallized with a matrix compound, and the mixture is placed on a probe tip and bombarded with a laser beam. The matrix is suitable for absorbing the laser wavelength, therefore, the energy from the laser pulse is absorbed by the matrix. Common matrices used commercially are a-cyano-4-hydroxycinnamic acid (CHCA), 3,5-dimethoxy-4-hydroxycinnamic acid (sinapic acid), and 2,5- dihydroxybenzoic acid (DHB).
Contd...
288 Bioanalytical Techniques The energy absorbed by the matrix is transferred to the sample so that it gets ionized and desorbs from the surface. The matrix plume carries the gas phase sample ions along with it to the detector which are analysed by the Time-of-Flight-(TOF) analysers. The TOF mass analyser separates ions in time as they travel through a flight tube. In the flight tubes, they are accelerated by applying an electric field into a field free drift region and are resolved according to their differential velocity depending upon the m/z value. The mass spectrum is obtained by measuring the detector signal for each pulse of ions produced in the source region, as a function of time (Figure B).
Figure B
9.5.3
Schematic representation of matrix-assisted laser desorption ionization TOF mass spectrometry
Interpretation of Mass Spectra
9.5.3.1 Analysis of peaks A mass spectrum is presented as a bar graph in which the relative abundance of ions taken on the ordinate (y-axis) is plotted against m/z taken on abscissa (x-axis). Most of the ions formed in a mass spectrometer have a single charge, so the m/z value is equivalent to mass itself. Modern mass spectrometers easily resolve ions differing even by a single atomic mass unit, thus providing completely accurate values for the molecular mass of a compound. (i) Molecular ion peak or parent peak The molecular ion is easily identified as the ion with highest m/z value unless it represents an impurity or an isotopic ion. Conventionally, it is used for calculating the molecular weight
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289
of a compound using the masses of the most abundant isotopes of all the elements present in it. Usually, the atomic mass is represented as the nominal mass in mass spectrometers which is obtained by rounding it off to the nearest digit. The parent peak represents the mass of the molecular ion produced by loss of an electron, that is, M+1 is often too weak to be seen in the high energy spectra that necessarily bring about fragmentation of the molecule. The intensity of the molecular ion is high for chemically stable molecule. Molecules containing p bonds tend to generate more stable ions, therefore, the relative heights of parent peaks in some classes of organic compounds can be presumed to decrease in the following order: Aromatics > Alicyclic > Ketones > Amines > Esters > Ethers > Carboxylic acids > Alcohols The nitrogen rule A molecular ion signal typically follows the nitrogen rule which states that if an uncharged organic molecule containing one or more atoms of the elements C, H, O, N, S, P, or a halogen has odd number of nitrogen atoms, it will have an odd-numbered nominal mass. If it contains even number of nitrogen atoms or no nitrogen atoms at all, then it will have an even-numbered nominal mass. See Figure 9.22. The ring rule It is a general rule applied to identify whether the molecular ion consists of even or odd number of electrons. If the molecular formula is given as Cx HyNzOn, then according to the ring rule, the number of rings plus double bonds is given by x – y/2 + z/2 + 1. This comes to be a whole number for an odd-electron ion, and a number ending in 1/2 for an even-electron ion.
Figure 9.22
Molecules depicting difference in nominal mass depending on the presence of nitrogen
290 Bioanalytical Techniques
.
(ii) Fragmentation and base peak Due to bombardment with high-energy electrons, molecules often undergo fragmentation and, thus, the parent peak is either weak in intensity or is completely missing from the spectrum. The molecular ions are energetically unstable and some of them cleave into smaller fragments. M+• Æ X+ + Y• The most intense peak in the spectrum is the base peak, which refers to the ion with the greatest abundance, assigned as 100 per cent. The heights of all other peaks are measured with respect to the base peak. Cleavage is more likely to occur at carbons exhibiting higher degree of branching and reflect the stability of the ion. Although the molecular ion is useful for identification, it usually does not provide any structural information about an unknown species. The structural information can be obtained from the fragmentation patterns of the mass spectrum. (iii) The isotope peaks Many elements have more than one naturally occurring stable isotopes which produce peaks in the mass spectrum other than the “main” peak, the one due to the most abundant isotope. Molecules having heavy isotopes give peaks corresponding to their m/z values that is, peaks occurring at (M+1) or (M+2). The intensities of these peaks are relative to their natural abundance. For example, the isotopes of chlorine 35Cl:37Cl occur in the abundance ratio 75.8:24.2 and their peaks give the intensity ratio 100:31.9 (Figure 9.23). Table 9.3 gives natural abundance of isotopes of some commonly occurring elements classified as M, M+1, and M+2 contributors depending upon their peak intensity in the MS spectra.
Chlorine
35
Cl
75.5
37
Cl
Chloro benzene + 35 Cl M · 112 ( Cl) + (M + 2) · 114 (37Cl) M+· : (M + 2)+· = 3:1
24.5
~ 3:1
112
114
m/z
Figure 9.23
Isotope peaks of chlorinein chlorobenzene
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291
Table 9.3 Natural abundance of isotopes of some elements M contributors
M+1 contributors
Isotope
Natural abundance Isotope (%)
1
99.9855
2
98.893
13
H
12 14
C
N
16
O
1.107
14
N
0.366
18
O
0.037
34
0.76
37
C
99.634
17
33
95.0
35
Cl
75.77
9.5.3.2
3
99.759
S
abundance Isotope
0.015
H
15
32
Natural (%)
M+2 contributors
S
Mass spectra of some compounds
(i) Methane (Figure 9.24) m/z
Ion/Fragment
17
Molecular ion (M+1)
16
Molecular ion (M); base peak
15
Fragment (M - H)
14
Fragment (M - 2H)
13
Fragment (M - 3H)
12
Fragment (M - 4H)
Figure 9.24
Mass spectrum of methane
H C O
Natural abundance (%) ppm ppm 0.204
S
4.22
Cl
24.23
292 Bioanalytical Techniques % Relative abundance
43
75
57
29
50
85 71
25
99
10
20
30
40
50
60
70
80
90
100
m/z
Figure 9.25
Mass spectrum of n-octane showing main peaks
(ii) n-octane: CH3CH2CH2CH2CH2CH2CH2CH3 (Figure 9.25) m/z 114 Æ M+ 99 Æ M - CH3 85 Æ M - CH2CH3 71 Æ M - CH2CH2CH3 57 Æ M - CH2CH2CH2CH3 43 Æ CH3CH2CH2+ 29 Æ CH3CH2+ (iii) Ethanol: CH3CH2OH (Figure 9.26) m/z 46 Æ M+ parent ion 45 Æ M - H 31 Æ M - CH3, CH2OH+ 29 Æ M - OH, CH3CH2+ 28 Æ M - H2O
114
110
120
Spectroscopy II
Figure 9.26
293
Mass spectrum of ethanol showing main peaks
(iv) Ethyl ether: CH3CH2OCH2CH3 (Figure 9.27) m/z 74 Æ 59 Æ 45 Æ 29 Æ
M+ M - CH3 M - CH2CH3 CH3CH2+
% Relative abundance
45
75
50 29 59
25
74
10
Figure 9.27
20
30
40
50
60 70 m/z
80
Mass spectrum of ethyl ether showing main peaks
90
100
110
120
294 Bioanalytical Techniques (v) Toluene: C7H8 (Figure 9.28) m/z 92 Æ M+ (C7H8+) 91 Æ M – H (C7H7+) % Relative abundance
91
92
80 60 40 20 0.0 15
Figure 9.28
30
45
60 m/z
75
Mass spectrum of toluene showing main peaks
90
105
10 Radioisotope Techniques
10.1
INTRODUCTION AND APPLICATIONS The discovery of radiation in 1896 by Henry Becquerel kindled an era of amazing scientific breakthroughs. Techniques that employed radioactivity illuminated the path of discoveries in all walks of science. Though the wartime application of this discovery is that of mass destruction, the peace time applications have served humanity by saving millions of lives. Few realize that radioactivity is used as a diagnostic tool and to treat deadly diseases such as cancer. The discovery of X-rays has enabled the visualization of skeletal bones for fractures. Today physicians are better equipped with several diagnostic tests based on radioimmunoassay and scanning tests for rapid and certain diagnosis. In fact the atomic structure, the basic unit of matter, could be unravelled by the use of radioactivity as a tool. Several new elements have been added to the periodic table by this discovery. History has witnessed at least 15 Nobel prizes awarded for pioneering discoveries that used radioactivity. In the first part of this chapter, we will discuss some of the brilliant applications of radioactivity.
10.1.1
Medical Uses
Radioisotopic techniques are extensively used in medicine so much so that a new branch of medicine known as nuclear medicine has been established. This branch of medicine uses radiation for diagnosis or treatment of diseases. (i) As diagnostic agents Radioactive isotopes can be detected in extremely small amounts. If a drug is labelled with a radioisotope, its passage through the body can be mapped and valuable medical information can be gathered about the functioning of organs. For example, 123I (iodine-123) is used in determining the state of health of the thyroid gland.
296 Bioanalytical Techniques (ii) Treatment of cancer Radiological damage caused by g rays is exploited in the treatment of cancer. If the radioactive element (60Co, 137Cs, or 99Tc) can be introduced into a tumour, its radiation will kill tumour cells. Moreover, since radiation damage is multiplied in rapidly growing cells such as tumour cells, healthy cells escape significant damage. For example, radio-cobalt (60Co) is used in the treatment of brain tumour, radio-phosphorous (32P) in bone diseases, and radio-iodine (131I) in thyroid cancer. Five Nobel laureates have been intimately involved with the use of radioactive tracers in medicine. (iii) In sterilization Bacteria and other disease-carrying organisms can be destroyed by irradiating them with g rays. The process is used for sterilizing medical instruments that would be damaged by heat sterilization. A portable source of g rays for sterilization is radio-cobalt (60Co). g rays are used for sterilizing medical and research equipment, as bacteria are killed by ionizing radiation.
10.1.2
Research in Biology
The great value of radiation in biology arises from the fact that radiation can be easily and accurately identified and measured in extremely minute amounts. Therefore, radioisotopes are used extensively in molecular biology as tracers to elucidate biological pathways and biomolecules. Radioisotopes have the same chemical properties as non-radioactive isotopes of the same element, and, thus, various compounds can be labelled with minute amounts of a radioisotope, which are easily detectable and quantified. They can be incorporated into DNA, RNA, and protein molecules, both in vivo and in vitro. As a consequence, the presence or metabolism of macromolecules under study can be “traced” and investigated.
10.1.3 Archaeology The age of archaeological or geological artefacts can now be determined by “carbon dating”—a remarkable technology that uses radioactive carbon. 14C (carbon-14) is present in small amounts in the atmosphere (as CO2), where it is produced by the irradiation of 14N by neutrons of cosmic rays. This dating technique is a valuable tool for determining the age of archaeological findings.
10.1.4 Agricultural Uses Radiations from radioisotopes are used for killing insect pests in food grains. Similarly, certain foodstuffs may be sterilized using this method that keeps the food fresh for an appreciable period of time. Radioisotopes can be used for determining the mechanism of action of fertilizers on plants. Moreover, they
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can be used for improvement of crops by different molecular techniques. Several high yielding varieties of pulses, oilseeds, rice, and jute have been indigenously developed using radioisotopic techniques.
10.1.5
Industrial Uses
Radiography of process equipment gives important information regarding wear and tear of the equipment parts. Radiotracers are used for measuring and monitoring sediment transport at ports, measurement of flow, and hydrology. Machinery used for making sheets such as paper, plastic, metals, employ methods based on radiography for measurement and control of sheet thickness. In industry, the tracer technique is used for testing the uniformity of mixtures. For example, for testing a food or pharmaceutical mixture, a small quantity of short-lived radioisotopes such as 24Na or 56Mn are added to the ingredients. Several different samples of the final products are then tested for radioactivity by means of a Geiger–Müeller counter. If each sample gives the same count, the mixture can be considered to be homogeneous. In the construction industry, the tracer technique is used for measuring the water level and location of rebar. After having discussed the remarkable applications of radioactivity, we will now understand the importance of bioanalytical techniques that employ radioactivity.
10.2
STRUCTURE OF ATOM AND RADIOACTIVITY Matter is composed of elements and all elements are composed of atoms. The atom can be thought of as a system comprising a positively charged nucleus and negatively charged electrons orbiting around it. The simplified Bohr’s model is composed of three fundamental particles called protons (positively charged), neutrons (neutral), and electrons (negatively charged). Each atom has the same number of protons as it has electrons. Thus, an atom is electrically neutral with the total positive charge in the nucleus equal to the total negative charge of its electrons. Also, the total number of electrons of an atom in its neutral or non-ionized state is equal to its atomic number. These electrons determine the chemical property of the atom. Each element has a unique number of protons and electrons. Each of these particles has a mass of approximately one atomic mass unit (amu) (1 amu = 1.66 × 10 –24 g). Many of the chemical elements have a number of isotopes. The isotope of an element has the same number of protons in their atoms (that is, atomic number) but different atomic masses due to different number of neutrons. For example, phosphorous has seven different isotopes. Each of the isotopes has 15 protons, while the number of neutrons in the isotopes varies from 28 to 34.
298 Bioanalytical Techniques There are 82 stable elements and about 275 stable isotopes of these elements. Many nuclides are unstable because the ratio of neutrons to protons produces a nuclear imbalance (too many protons or too many neutrons in the nucleus). These unstable isotopes attempt to become stable by rearranging the number of protons and neutrons in the nucleus and achieving a more stable ratio. The excess energy is ejected from the nucleus as radiation. This process is known as radioactive decay and the isotope undergoing radioactive decay is said to be radioactive isotope.
10.2.1
Natural and Artificial Radioisotopes
Numerous radioisotopes occur naturally in the atmosphere. Carbon, for example, occurs as Carbon-12 (12C) with six protons and six neutrons and is completely stable. Interaction with cosmic rays in the atmosphere produces 14C, an isotope of carbon that consists of six protons and eight neutrons. 14C decays into 14N which consists of seven protons and seven neutrons by changing neutron to a proton and emitting a b particle. Measurement of the relative abundance of these two isotopes of carbon is the basis of carbon dating. There are a number of unstable natural isotopes arising from the decay of primordial uranium and thorium. Each of these series ends with a stable nuclide of lead. Figure 10.1 shows the decay series from Uranium-238. When a novel combination of neutrons and protons is produced artificially, the atom is unstable and is called a radioactive isotope or radioisotope. Radioisotopes can be manufactured in several ways. The most common method is by neutron activation in a nuclear reactor. This involves the capture of a neutron by an atom’s nucleus resulting in excess of neutrons in the nucleus. Some radioisotopes are manufactured in a cyclotron in which protons are introduced to the nucleus resulting in deficiency of neutrons. There are presently 1800 known natural and artificial radioisotopes. About 200 artificially produced radioisotopes are used on a regular basis.
10.3 TYPES OF RADIOACTIVE DECAY The nucleus of a radioisotope usually becomes stable by emitting a particle, b particle, or positron or there may be a combination of the decay particles. These particles may be accompanied by the emission of energy in the form of electromagnetic radiation known as g rays. This process of nuclear disintegration can be one of the four different types: a radiation, b radiation, electron capture, or g radiation.
Radioisotope Techniques
Figure 10.1
10.3.1
299
Natural chain of radioactive decay
a radiation
Alpha radiation is particulate radiation with a very large mass (atomic mass of 4). a decay occurs in elements with high atomic mass, those with a heavy nucleus, and those with very high neutron to proton ratio; for example, radon or uranium. The helium nucleus (42He) is a massive particle of 4 amu with a +2 charge. After a radiation, the atomic mass number (as denoted by A in Equation 10.1) decreases by 4, the number of protons decreases by 2, and the number of neutrons (N) decreases to N – 2. A X Z N
Parent isotope
Æ
4 2+ He 2 + 2 a particle
A–4 Y Z–2 N – 2
10.1
Daughter isotope
Due to greater mass of a particles, this radiation can impart large energy to matter though which it travels. When a particles intercept matter, the energy
300 Bioanalytical Techniques loss from the particle occurs because the particle incurs elastic collisions with atomic electrons of the target material. This usually leads to ionization of the target, ultimately culminating in the a particle picking up the two electrons. Because of their biologically damaging characteristics, a emitting isotopes have little use in nuclear pharmacy. Moreover, a particles are difficult to detect because of their short range, so they have little utility for imaging. They are absorbed by 10 cm of air, 0.01 mm lead, or a sheet of paper.
10.3.2
b radiation −
In b decay (negatron decay), the weak nuclear interaction converts a neutron into a proton and in the process emits an electron and an anti-neutrino (Ùe). Therefore, as a result of positron emission the nucleus gains a proton, the atomic number increases by 1, and the atomic mass number remains constant. b particle has only about 1/8000 of the mass of an a particle and only half of the charge. Its penetrating power is, therefore, higher than a particle (absorbed by 1 m air, 0.1 mm lead, or 3 mm aluminium sheet) whereas they are less ionizing. __
n Æ P + 0 –1e – + ne
10.2
In b+ decay (positron decay), a proton is converted into a neutron, a positron, and a neutrino. P Æ n + 0 +1e + + ne
10.3
If the proton and neutron are part of an atomic nucleus, the decay process transmutes one chemical element into another. For example, Cesium-137 has 55 protons and 82 neutrons in its nucleus. A nucleus of 137Cs is unstable and decays to a lower energy nucleus, Barium Ba-137, which has 56 protons and 81 neutrons. One of the neutrons changes into a proton and emits an electron and an antineutrino. __ 137 137 Cs Æ Ba + 0 –1e – + ne (b – that is beta emission) 55 56
10.4
As a result of positron emission, the nucleus loses a proton and gains a neutron, and the mass number remains the same. An example of such an isotope decaying by positron emission is 22Na which decays into neon-22 with the emission of a positron. 22 Na Æ 22 Ne + 0 +1e + (b+ that is positron emission) 11 10
10.5
Positron emitters are detected by the same instruments used for detecting g radiation. They are used in biological sciences to spectacular effect in brain
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301
scanning with positron emission tomography (PET scanning) technique used for identifying active and inactive areas of the brain.
10.3.3
Electron Capture
In this form of decay, a proton captures an electron orbiting in the innermost K shell. Proton + Electron = Neutron + X-ray
10.6
The proton becomes a neutron and electromagnetic radiation (X-rays) occurs. For example, Iodine-125 decays by electron capture to an excited state of Tellurium-125. 125 I Æ 125 Te + X-ray 53 52
10.3.4
10.7
g Radiation
Gamma rays are electromagnetic waves emitted from nuclear reactions. These waves travel at the speed of light and have high penetrating power. Unlike aand b radiations, g rays are not particulate in nature. These waves can pass through the interstitial spaces of mammalian tissue without causing disruption. However, in passing through sufficient thickness of living tissue the possibility of interaction and, thus, ionization greatly increases. The g radiation has low ionizing power but high penetration. For example, the radiation from 60Co can penetrate 15 cm of steel. On decaying, Iodine-131 most often expends decay energy by transforming into the stable Xenon-131 in two steps, with gamma decay following rapidly after beta decay. 131 I Æ 131 Xe + b– + g 53 54
10.4
10.8
INTERACTION OF RADIOACTIVITY WITH MATTER The energetically charged particles such as a particles and b particles, while passing through matter, lose their energy by interacting with the orbital electrons of the atoms in the matter. During this process, the atoms are ionized wherein the electron in the encounter (that is hit by the radiation) is ejected or raised to a higher energy state resulting in the excitation of the atoms. In both excitation and ionization processes, chemical bonds present in the molecules of the matter may be ruptured, forming a variety of chemical entities. The lighter charged particles (for example, b particles) follow a zigzag path in the matter, whereas the heavier particles (for example, a particles) move in a straight path, because of its heavy mass and charge. A positron loses its energy by interaction with
302 Bioanalytical Techniques electrons of the absorber atoms and comes to almost rest by combining with an electron of an absorber atom. At this instant, both particles (b+ and e−) are annihilated as a result of matter–antimatter encounter to produce two photons which are emitted in opposite directions. (i) a particles Alpha particles carry considerable energy of around 3–8 MeV. They react with matter in two ways. (a) They cause excitation where the energy from a particle is transferred to orbital electrons of the interacting atoms resulting in elevation of these electrons to higher orbitals. These electrons fall back to their original paths emitting energy as photons. (b) They ionize the atoms in their path by removing the electrons from the atom and leading to the formation of ions. Since a particles are heavy, they are capable of causing intense ionization and excitation. However, they are not very penetrating as their paths are obstructed by the atoms of the matter from which they travel. (ii) b particles These are light and rapidly moving particles that carry a single negative charge. These particles, also result in excitation and ionization of matter in which they travel. However, since they are much smaller as compared to a particles, they are less likely to encounter the atoms of matter in their paths. Therefore, they are less ionizing. However, they are more penetrating. The energy of emitted particles is different for different isotopes. (iii) Bremsstrahlung radiation When high atomic number materials absorb high-energy b particles, the absorber gives out a secondary radiation, an important phenomenon for X-ray generation, called Bremsstrahlung radiation. (iv) g rays They bear no charge or mass; therefore, they cause considerably less ionization as compared to particulate radiation. g rays interact with matter to create secondary electrons that behave as per negatron emission. In the spectrum of electromagnetic radiations, g radiations are high-frequency radiations and interact with matter by three mechanisms: photoelectric, compton, and pair production. (v) Photoelectric process In this process, a g radiation, while passing through an absorber, transfers its entire energy primarily to an inner shell electron (for example, the K-shell) of an absorber atom and ejects the electron. The vacancy in the shell is filled in by the transition of an electron from the upper shell, which is followed by emission of the energy difference between the two shells as characteristic X-rays. (vi) Compton scattering process In a compton scattering process, a g radiation with somewhat higher energy interacts with an outer shell electron of the absorber atom transferring only part of its energy to the electron and
Radioisotope Techniques
303
ejecting it. The ejected electron is called the compton electron and carries a part of the g ray energy minus its binding energy in the shell. As the energy of the g radiation increases, the photoelectric process decreases and the compton scattering process increases but the latter also decreases with photon energy above 1.0 MeV or so. (vii) Pair production When the g-ray energy is higher than 1.022 MeV, the photon interacts with the nucleus of an absorber atom during its passage through it and produces a positron and an electron. This is called pair production. The excess energy beyond 1.022 MeV is shared as kinetic energy between the two particles. (viii) Attenuation of g radiations When g radiations pass through the absorber medium, they undergo one or a combination of the above three processes (photoelectric, compton, and pair production) depending on their energy, or they are transmitted out of the absorber without any interaction. The combined effect of the three processes is called the attenuation of the g radiations.
10.4.1
Biological Effects of Radiation
Radiation is potentially hazardous to the living system. We, therefore, stress the need for best practices for the safety of personnel involved in research that deals with radiation as well as those who are in risk of occupational exposure. The regulatory bodies all over the world base their policies on the principle that all forms of radiation, even very minute measure, may be capable of producing stochastic detrimental effects. Ionizing radiation absorbed by biomolecules leads to removal of electrons from the atoms resulting many a times in breaking the molecular bonds and, hence, damaging the tissue. It is not certain that ionizing radiation would result in damaging a vital part of the cell. However, it has been proved to be the cause of damage to deoxyribonucleic acid, the biomolecule that encodes the genetic instructions used in the development and functioning of living organisms. However, all the cells are not equally sensitive to damage caused by radiation. In general, cells which divide rapidly and/or are relatively nonspecialized tend to show effects at lower doses of radiation then those which are less rapidly dividing and more specialized. Examples of the more sensitive cells are haematopoietic system that produces blood, the most sensitive biological indicator of radiation exposure.
10.5
KINETICS OF RADIOACTIVE DECAY Radioactive decay is a spontaneous process and it occurs at a rate characteristic of the source, defined by the rate constant (l), the fraction of an isotope
304 Bioanalytical Techniques decaying in unit time, t –1. It follows first-order reaction kinetics. Unlike many chemical reactions, the rate of radioactive decay is not affected by temperature or pressure. A given sample of any radioisotope, the number of decay events (− dN) expected to occur in a small interval of time (dt) is proportional to the number of atoms present (N), that is –
dN ____ μ N dt
10.9
As explained previously, the decay rate is a characteristic of the radionuclides decay; therefore, each has its value of l. –
dN ____ = ldt N
10.10
The negative sign indicates that N decreases as time increases, as each decay event follows one after another. The solution to this first-order differential equation is the function N(t) = N0e –lt = N0e –t/T
10.11
where N0 is the value of N at time t = 0. The number of atoms disintegrating at any time is proportional to the number of atoms of the isotope N present at time t. Since the number of atoms N is always decreasing, the rate of decay, that is, disintegrations per minute (dpm), also decreases with time. The profile of graph obtained for number of radioactive atoms against is an exponential decay curve (Figure 10.2). Equation 10.3 can be written in a logarithmic form as follows In Nt /N0 – lt
10.12
A plot of natural logarithm of the number of dpm against time would be a straight line and a negative slope (–l). In general practice, radioactive decay constant is expressed as the half-life (t½). This is defined as the time taken for the activity to fall half of the original value. When Nt in Equation 10.12 is equal to one-half of N0, then t will be equal to the half-life of the isotope. Thus In 1/2 = lt1/2
10.13
or t1/2 = 0.693l
10.14
Radioisotope Techniques
Figure 10.2
305
Exponential decay of radioisotopes
The half-life of a radioisotope can range from a small fraction of a second to millions of years (Table 10.1). Since the half-life of a particular radionuclide is unique to that radionuclide, the data of the half-life may assist in identification of the radionuclide. Table 10.1 Common radioactive isotopes produced during nuclear reactions Isotope
Half-life
Isotope
Half-life
Isotope
Half-life
Relatively short half-lives Strontium-89
54 days Zirconium-95
65 days
Ruthenium-103
40 days Rhodium-103
57 minutes
Iodine-131 Barium-140
8 days Xenon-133 13 days Lanthanum-140
8 days 40 hours
Niobium-95
39 days
Rhodium-106
30 seconds
Tellurium-134
42 minutes
Cerium-141
32 days
Year to century-scale half-lives Hydrogen-3 Ruthenium-106 Promethium-147 Curium-224
12 years Krypton-85
10 years
Strontium-90
29 years
1 year Cesium-137
30 years
Cerium-144
1.3 years
Americium-241
440 years
2.3 years Plutonium-238
85.3 years
17.4 years Longer half-lives
Technecium-99 Plutonium-240
6
2 × 10 years Iodine-129 6,500 years Americium-243
1.7 × 107 years 7,300 years
Plutonium-239
24,000 years
306 Bioanalytical Techniques 10.6
UNITS OF RADIOACTIVITY
10.6.1
System of Units of Radioactivity (Becqueral)
The international system of units (SI) specifies Becquerel (Bq) as the unit of radioactivity. 1 Becquerel is defined as 1 radioactive disintegration per second (dps). Radioactivity was previously measured in curie (Ci) as defined as 1 g of radium-226 or 3.7 × 1010 dps. 1 Becquerel = 1 radioactive disintegration per second = 2.703 × 10 –11 Ci
10.6.2
Radiation Exposure (Roentgen)
Exposure describes the amount of radiation travelling through air. Many radiation monitors measure the exposure given in the units of Roentgen (R) or coulomb/kilogram (C/kg). It is the amount of g or X-rays required for producing ions resulting a charge of 0.000258 C/kg of air under standard conditions. This unit has been named after Wilhelm Roentgen, the German scientist who discovered X-rays in 1895. Radiation dose is the magnitude of radiation exposure. It has two categories: absorbed dose and biological dose.
10.6.2.1 Absorbed dose (Gray) It is defined by the amount of energy deposited in a unit mass of exposed tissue. Absorbed dose is measured in gray (Gy), that is equivalent to 1 J/kg or rad, which is 100 erg/g.
10.6.2.2
Biological dose or dose equivalent (REM)
This dose takes into account that the biological damage due to radiation not only depends on the total energy deposited, but also on the rate of energy loss per unit distance travelled. For example, a particles do much more damage per unit energy deposited than electrons. For a and g radiations, the dose equivalent is the same as the absorbed dose. This effect has been defined by a quality factor Q. Over a wide range of incident energies, Q is taken to be 1.0 for electrons [and for X-rays and g rays, and 20 for a particles (Table 10.2)]. For neutrons, the adopted quality factor varies from 5 to 10, depending on neutron energy. The biological impact is specified by the dose equivalent H, which is the product of the absorbed dose D and the quality factor Q. H = Q × D
10.15
The SI unit of biological dose or dose equivalent is expressed in terms of Sievert (Sv) or in Roentgen equivalent man (rem).
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307
Table 10.2 Units of radioactivity Type of radiation
Rad
Q factor
rem
X-ray
1
1
1
g ray
1
1
1
b particles
1
1
1
Thermal neutrons
1
5
5
Fast neutrons
1
10
10
a particles
1
20
20
1 Sv = 100 rem
10.16
Biological dose equivalents are commonly measured in 1/1000th of a rem (known as a millirem or mrem). However, dose equivalent is still not the whole story. If only a part of the body is irradiated, the dose must be discounted with an appropriate weighting factor if it is to reflect the overall risk. The discounted dose is termed the effective dose equivalent or just the effective dose, expressed in rem or Sv. The mode of interaction of particulate and non-particulate radiation with matter is different.
10.7
DETECTION AND MEASUREMENT OF RADIOACTIVITY
10.7.1
Ionization Chambers
Alpha, beta, and gamma radiations fall into the category of ionizing radiations. Ionizing radiation creates ion-pairs in the matter as it passes through. In 1908, Hans Geiger developed a device for the detection of a particles based on ionization of matter due to interaction with ionizing radiations. Later by 1928, Walther Müeller, a student of Geiger, improvised the counter to the modern Geiger–Müeller counter that enabled the detection of any kind of ionizing radiation. There are three different types of ionization detectors: (i) ionization chambers, (ii) proportional counters, and (iii) Geiger–Müeller counters. In all of the above types, a volume of gas is surrounded by one electrode and contains in its interior a second electrode. A potential difference is maintained between the electrodes, and the current needed to keep that potential constant is measured by electronic circuits. The ion pairs formed due to incident radiation are in random thermal motion and diffuse in the chamber. The free electrons collide with ions and neutral gas molecules resulting in either charge transfer or recombination of the ions.
308 Bioanalytical Techniques Charge transfer occurs when a positive ion encounters a neutral molecule resulting in transfer of electrons from the neutral molecule to ion. Whereas, when a free electron collides with a gas molecule that attaches electrons easily, a negative ion is formed. The released electrons and ions are collected in an external electrical field on two electrodes mounted in the gas. This produces an electrical signal at the electrodes, which can either be measured as current in the case of a high count rate or as electrical voltage pulse of certain pulse height V. The three types of ionization chamber detectors differ only in the magnitude of applied voltage across the electric field in the chamber. Pulse height recorded is directly proportional to the number of ion pairs produced by the ionization event. This can be adequately explained using the relation between the pulse height versus applied voltage as given in Figure 10.3. The curve that is obtained can be divided into distinct regions: recombination region occurs at very low voltages in which the applied voltage is low and the ions move slowly towards the electrodes and tend to recombine to form neutral atoms or molecules. Therefore, the pulse height is less than it would have been if all the ions originally formed reached the electrodes. Gas ionization instruments are, therefore, not operated in this region of response. In the second region of the curve, known as the ionization region, the voltage is sufficient to minimize recombination and collect the ion pairs on the electrodes. This is the region in which ionization chambers operate. The signal amplitude generated is extremely small in pulse mode which requires a sophisticated low-noise preamplifier and pulse processing electronics. If the incident radiation is fast electron or g ray, then the pulse amplitude is even smaller. If the voltage is further increased, electron gains a very high velocity and this velocity is sufficient to induce ionization of other atoms or molecules. This phenomenon is referred to as Townsend avalanche and this region of the curve is called the proportional region since the gas amplification factor is proportional to the applied voltage. This is the principle behind proportional counters that operate in the proportional region as shown in Figure 10.3. Amplification eliminates the need for very sensitive current detecting circuits. The increase in the applied voltage results in non-linear effects (limited proportional region). The free electrons are quickly collected due to high mobility while positive ions are slowly moving, which takes a long time. This results in non-linearity and, therefore, this region is not used for ionization chambers. If the voltage is further increased, the field strength is so great that the discharge, when initiated, continues to spread until further amplification cannot
Radioisotope Techniques
Figure 10.3
309
Operational regions of gas filled radiation detectors
occur due to accumulation of thick positive ion sheath surrounding anode that results in decrease in electric field (Figure 10.5). Each pulse is of the same amplitude. This is the region in which Geiger–Müller are operated but since the pulse amplitude does not reflect the energy absorbed by the detector by each interaction, the Geiger–Müller pulses contain no energy or particle information. I n t he cont i nous d isch a r ge reg ion, t he appl ied volt age is so high that once ionization occurs in the gas, there is a continuous discharge of electricity. Therefore, this region too is not useful for the detection of radiation.
10.7.1.1
Design and working of Geiger–Müeller counter
A Geiger–Müller tube has an anode and a cathode housed in a tube-like shell that contains a noble gas such as argon at low pressure (50–100 torr). The intense electric field near the anode collects the electrons to the anode and repels the positive ions (Figure 10.4). Electrons reach the anode within a fraction of seconds (µs) whereas the positive ions that are heavier travel more slowly (reach cathode within milliseconds, ms) forming a sheath of massive positive ions surrounding the anode. The temporary build-up of a positive charge surrounding the central anode terminates production of additional avalanches by reducing the field gradient near the centre wire below the avalanche threshold. (i) Quenching If the ions reach the cathode with sufficient energy, they can liberate new electrons, starting the process all over again, producing an endless continuous discharge that would render the detector useless. Therefore, a “quenching material” is used in small amount (5–10 per cent)
310 Bioanalytical Techniques
Figure 10.4
Schematic of Geiger– Müeller counter
Figure 10.5
Formation of avalanche in a Geiger– Müeller counter
along with the fill gas in the ionization chamber. The quenching material has a lower ionization potential than that of the fill gas. When the positive ions encounter the quench gas molecules, they tend to transfer their charge to the quench gas molecules. The quench molecules are more complex in nature and the excess energy released during neutralization is used for their degradation rather than for liberation of a free electron from the cathode surface.
Radioisotope Techniques
311
(ii) Dead time The dead time of a Geiger–Müller tube is defined as the period between the initial pulse and the time at which a second Geiger discharge can be developed. In most Geiger–Müller tubes, it is of the order of 50–100 µs. The detector is unresponsive while the ion sheath migrates outwards for recovering the field gradient above the avalanche threshold and starts the process all over again. The Geiger–Müller counters are regularly applied in hospitals and industries for monitoring radiation. This is because they offer benefits such as simplicity of operation, economy, and detection of all three types of radiation. They do not require highly stabilized electric circuit as with the ionization detectors. However, the main disadvantage is that they have a considerable dead time of 200–300 µs. Also they cannot be used for determining the type of radiation as the pulse generated is independent of the energy of incident radiation.
10.7.2
Scintillation Counters
Scintillators are substances that emit light when colliding with an ionizing particle. Scintillation-based radiation counters are the most sensitive and versatile instruments for the detection and quantification of radioactivity. These are now standard laboratory devices in the life-sciences establishments for measuring radiation from g and b emitting nuclides. Scintillation detectors have high efficiencies of detection, easier sample preparation, amenability to automation including data processing, and ability to detect different nucleotides. These advantages have made this technique very popular. Modern scintillation counters are of two types: solid scintillation counter (SSC) and liquid scintillation counter (LSC) detectors depending upon the scintillation material used. In SSC, scintillation material comprises the crystals of inorganic material such as alkali halide crystals, mainly sodium iodide. The high atomic number and density of certain inorganic crystals make them ideal for detecting g radiation. On the other hand, LSC has an organic liquid or a plastic as the scintillation material and this is preferred for the detection of b particles and neutrons. Scintillation material is doped with activators that increase the efficiency of generation of photons.
10.7.2.1
Solid scintillation counters
In solid scintillation counters, the phenomenon of scintillation occurs due to the nature of electronic band structure in the crystals. The scintillation material absorbs energy of incident radiation to excite the electrons from valence electronic energy band to the conduction band (Figure 10.6). This
312 Bioanalytical Techniques Conduction band (empty)
e Activator states
Excitation Photon
Valence band (full)
Hole (Energy gap about 4 eV)
Figure 10.6
Energy bands in scintillation material
leaves a hole in the valence band. The electron-hole pairs are called excitons and are captured as a whole by the activator molecules leading to excitation of activator molecules. The excited activator molecules return to their ground state rapidly by emitting scintillation. This is the fast component of scintillation. The activator molecules are also called dopants or doping material. These are added to the scintillation material in trace amounts so that the light emitted is in the visible range (or range where photomultipliers are active). Some activator molecules achieve metastable energy states after capturing excitons. The delayed de-excitation of these metastable molecules results in light emission known as the slow component of scintillation. Copper, silver, cadmium sulfate, europium, cerium, and thallium are commonly used as activators or doping materials. A scintillation counter apparatus consists of a scintillator, a photomultiplier tube (PMT), an amplifier, and a multichannel analyser (Figure 10.7). Photomultipliers consist of vacuum glass tubes which contain a photocathode, numerous dynodes, and an anode. The photons emitted by the scintillating material strike the photocathode resulting in emission of electrons due to photoelectric effect. These electrons are directed to a number of electrodes called dynodes that are arranged in the PMT according to increasing positive potential (Figure 10.8). The electrons strike the first dynode to release secondary electrons. These electrons are attracted to second dynode and strike to release more electrons. This amplification continues through 10–12 stages. Due to this amplification, a voltage pulse can be generated across the external resistors.
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Photomultiplier tube
Sodium-iodide crystal of liquid scintillation fluid
Photocathode Optical window
Radiation
Figure 10.7
–+ Anode +HV
Light Photo- Dynode electron
time
Measuring device
v
Design of a solid scintillation counter
The magnitude of the pulse is directly proportional to the photon intensity emitted by the scintillators. This voltage pulse is amplified and recorded by a multichannel (counting windows) that classifies each voltage pulse. Each channel corresponds to a specific range of energy channels and counts with energies above or below set limits are excluded from a particular channel. Pulses are gathered into channels and the counts per minute (CPM) in each channel and energy distribution or spectrum is recorded. CPM is proportional to the amount of isotope in the sample, and the spectrum indicates the identity of the isotope. Solid scintillation can be applied for g radiation because it is highly penetrating and interacts with the sodium iodide–thallium (NaI–Tl) detector by photoelectric effect. Since the crystal must be protected from contamination by the sample, there is a need of barrier to prevent the contamination. However, Photomultiplier
e
Photocathode
Figure 10.8
Dynodes
Signal amplification in photomultipiler tube
Anode
314 Bioanalytical Techniques extremely low-penetrating power of a and b radiation limits the distance that can be traversed to reach the scintillators. Liquid scintillation counters can be used for these types of radiation.
10.7.2.2
Liquid scintillation counters
In liquid scintillation counters (LSC), the samples are dissolved or suspended in a solution of organic liquid or cocktail scintillators. This assures close contact between the radioactive particles and the scintillators. Liquid scintillation cocktails absorb the energy emitted by radioisotopes and re-emit it as flashes of light. Initially the solvent molecules absorb the energy of radiation and then transfer it to scintillators or phosphors that convert the absorbed energy to light. (i) The role of the solvent The solvent portion of an LSC cocktail comprises 60–99 per cent of the total solution. Solvent acts as an efficient collector of energy and conducts this energy to the phosphor molecules. An ideal solvent dissolves the phosphor to produce a stable solution and has minimum quenching effect. Quenching is the loss of counts due to sample or cocktail characteristics. Quenching may result from a variety of components in a sample. The typical LSC solvent is toluene. The p cloud of the toluene ring (or any aromatic ring) provides a target for b-interaction, which captures the energy of the incident particle. Thus, a b particle passing through a toluene solution leaves a trail of a number of energized toluene molecules. The energy from these molecules passes back and forth among the solvent ring systems, allowing efficient capture by dissolved phosphors. The solvents that are employed are aromatic organic compounds such as toluene, pseudocumene and phenyl-o-xylylethane (PXE), and possess aromatic rings to absorb the energy of incident radiation. (ii) The role of scintillators Primary scintillators, also known as flours or phosphors, are added at 0.3–1 per cent of the solution volume and they convert the captured energy into light. Common f luors include polyphenyl hydrocarbons, oxazole and oxadiazole aryls, PBD 2-Phenyl-5-(4-biphenylyl)-1,3,4-oxadiazole, butyl-PBD 2-(4-Biphenylyl)5-(4-tert-butylphenyl)-1,3,4-oxadiazole, PPO 2,5-diphenyloxazole. Secondary scintillators or activators or dopants are also included in the LSC cocktail. An activator captures florescence energy of the excited scintillator, and re-emits it as a longer wavelength signal. This increases the efficiency of scintillation and facilitates its detection. Common dopants or secondary scintillators used are: POPOP 1,4-bis(5-phenyl-2-oxazolyl) benzene, dimethyl-POPOP 1,4-bis(4-methyl-5-phenyl-2-oxazolyl)benzene.
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(iii) Quenching Chemical quenchers reduce the number of photons generated by radioactive particles. Quenchers can be classified as chemical quenchers or colour quenchers. Chemical quenchers absorb radioactive energy before it is converted to light. Colour quenchers absorb light in the range of the wavelength emitted by the scintillator. In this case the number of photons emitted is not changed, but the number reaching the photomultiplier tube is reduced. This leads to an underestimation of the total counts, and thus of the isotope present. It also leads to an apparent shift in the energy spectrum of the sample. (iv) Quench correction Various methods are available for quench correction. The most straightforward, but most laborious, is the use of an internal standard. A known amount of radioactivity, added to an unknown sample, will increase the dpm by a predictable amount. The difference between the increase in dpm observed and that expected is due to quenching, and allows the determination of counting efficiency for that sample. The drawback to the use of internal standards is that each sample must be counted twice. It is also inconvenient to add an internal standard to many vials. Many scintillation counters offer the use of an external standard to correct quenching. After initial counting, a strong g-emission source is placed next to the vial and the sample is counted again. The g rays cause secondary emission of compton electrons, which scintillate with the sample. The counts due to sample radioactivity are subtracted, leaving only the compton electron counts. The theoretical energy distribution of the compton electrons is compared with the measured energy spectrum to determine the extent of quenching. The samples must still be counted twice, but nothing needs be added to the vials, and the process may be carried out automatically by the counter. (v) Signal analysis The number of photons generated is directly proportional to the path length of the b particle, which is in turn determined by its emission energy. The intensity of each light pulse corresponds to the energy emission and the number of pulses per second corresponds to the number of radioactive emissions. Therefore, CPM is directly proportional to the amount of isotope present in the sample, and the spectrum indicates the identity of the isotope. The scintillation counter has pulse height analyser that classifies each pulse according to the number of photons in the pulse. The number of photons in a pulse corresponds to the energy of individual b-emission event. The user can conveniently set a pulse height analyser to reject all pulses of energy below X (where X is the threshold value) and to reject all pulses of energy above Y (where Y is the window under analysis). The Analyser set with a threshold and window for a particular
316 Bioanalytical Techniques isotope is known as a channel. Each channel corresponds to a specific range of b-energies (channels are also known as counting windows). The principle of the method is illustrated in Figure 10.9, where it can be seen that the spectra of two isotopes (referred to as S and T) overlap only slightly. Modern counters operate with multichannel analysers that record the entire energy spectrum simultaneously. Examples of pairs of isotopes that can be counted together are: 3H and 14C, 3H and 35S, 3H and 32P, 14C and 32P, and 35S and 32P. Dual-label counting has proved to be useful in many aspects of molecular biology (for example, nucleic acid hybridization and transcription), metabolism (for example, steroid synthesis), and drug development. (vi) Chemiluminescence and static electricity Chemiluminescence is caused by any chemical reaction that generates an excited product molecule, which decays to emit light. These reactions generate only a single photon, which may be quenched, or may reach the counter to register as a low-energy emission event. Spurious counts due to chemiluminescence can be detected by counting the samples twice with a period of about an hour interval between the counts. The rate of photon emission decreases drastically as chemiluminescent reaction proceeds. In contrast, even a short-lived isotope such as 32P will decrease its emissions only slightly over 24 hours. Phospholuminescence results from pigments in the sample absorbing light and re-emitting it; the solution is to keep the samples in the dark prior to counting.
Figure 10.9
Counting of dual-labelled samples
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Another source of false counts is static electricity. The energy from a static electric build-up can be released as a burst of light from the cocktail. Static can be minimized by wiping the vials with a wet paper towel (water dissipates the static electric charge) or by wiping with an antistatic laundry dryer sheet. (vii) Sample assay Liquid scintillation counting of discrete samples is straightforward. A sample is mixed with an appropriate volume of scintillation cocktail, and the mixture is placed in an LSC vial and counted. The best counting efficiencies are achieved when samples uniformly disperse into the cocktail to produce a clear, colourless emulsion of neutral pH. The scintillation counting can also be used for analysing the thin layer chromatography (TLC) plates. In a typical TLC experiment, the radioactivity is detected at two points: after TLC it is analysed by autoradiography and to locate radioactive spots. These spots are then scraped off of the plate and counted to provide quantitative information. (viii) Applications of scintillation counters Scintillation counters are used for measuring radiation in a variety of applications. • Hand-held radiation survey meters • Personnel and environmental monitoring for radioactive contamination • Medical imaging • National and homeland security • Border security • Nuclear plant safety • Radon levels in water • Oil well logging Several products have been introduced in the market utilizing scintillation counters for detecting potentially dangerous g-emitting materials during transport. These include scintillation counters designed for freight terminals, border security, and ports, weighbridge applications, scrap metal yards, and contamination monitoring of nuclear waste. Scintillation counting is a popular technique in biological work since it is very fast with practically no dead time. The counting efficiencies are relatively high ranging from about 50–90 per cent for low-energy to highenergy b-emitters. Samples in solid, liquid, suspensions, and gel forms can be easily detected using scintillation counters. Hundreds of samples can be counted and analysed automatically. However, scintillation apparatus is not only expensive but also higher cost per sample is incurred. Moreover, as the sensitivity is higher, a higher degree of noise is observed.
318 Bioanalytical Techniques (ix) Cerenkov counting Cerenkov counting is a relatively simple and cheap method for detecting b particles that can pass through a substance with a speed higher than that of light passing through the same substance. When a b-emitter has decay energy higher than 0.5 MeV, it causes water to emit a bluish white light referred to as Cerenkov light. It is possible to detect this light using standard liquid scintillation counters. Since there is no requirement for organic solvents and fluors, this technique is relatively cheap, sample preparation is very easy, and does not suffer from problems such as chemical quenching. Isotopes such as 32P can be counted using Cerenkov counting method. (x) Scintillation proximity assay The scintillation proximity assay (SPA) beads are designed in such a way that they can selectively bind the molecule under investigation with a scintillant. This is particularly useful with some types of radiation that do not travel long distances; for example, weak-energy emitters 3H and 14C. In usual LSC, if molecules containing such radioisotopes are in solution they do not result in scintillation as most of the energy is absorbed by the solvent. Therefore, these cannot be detected efficiently by this method. In turn when these are bound to the scintillant they can be detected easily. There are many applications of this technology such as studies involving enzyme assays and receptor binding.
10.7.3
Methods-based Upon Exposure of Photographic Emulsions “Autoradiography”
Ionizing radiation acts upon a photographic emulsion or film to produce an image similar to the one formed with visible light. The photographic emulsion or film contains silver halide crystals. The energy of incident radiation results in reduction of silver halide molecules to metallic silver forming a latent image. After treatment with suitable developers, silver grains can be visualized in the form of blackening of the photographic film. This technique is known as autoradiography. It is a very sensitive technique and has been used in a wide variety of biological experiments. For example, autoradiography of nucleic acids separated by agarose gel electrophoresis. Generally, weak b-emitting isotopes (for, example 3H, 14C, and 35S) are most suitable for autoradiography, particularly for cell and tissue localization experiments. A typical autoradiographic technique involves enclosing the gel or blot with X-ray film for at least 24 hours within an X-ray cassette. Development of the film reveals the gel or blot image. Note the use of radioactive or phosphorescent dots in recording the orientation of the film. (i) Choice of emulsion and film Autoradiography emulsions are solutions of silver halide that can be made to set solid by the inclusion of materials such as gelatine. This can be used for autoradiography of microscope slides.
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The shades of grey in the image are related to a combination of levels of radiation and length of exposure until a black or nearly black image results. Therefore, quantitative results can be obtained. However for successful application of this technique isotopes with energy of radiation equal to or higher than 14C are required. Direct autoradiography is ideally suited for recording emissions from 14C, 35 S, or 33P which emit b-particles with high energies when the b particles are not absorbed internally by the sample. However, most of the biological samples with low-energy emissions such as those from 3H, 35S, and 14C are internally absorbed within the sample, failing to reach the film. This problem is greatest for the very weak emissions of 3H but it also decreases the autoradiographic efficiency for 35S, 14C, or 33P when these isotopes are embedded within the thickness of a sample such as a polyacrylamide gel. These problems can be resolved using the technique of fluorography. (ii) Fluorography The detection of low-energy isotopes that cannot be usually detected by direct autoradiography can be made possible by fluorography.
Figure 10.10
Development of autoradiographic film
320 Bioanalytical Techniques A f luor, for example polyphenyl oxides (PPOs) or sodium silicate [Na2 (SiO2) nO] can be used for enhancing the image. The b particles excite the fluor molecules and emit light. This has been used for detecting radioactive nucleic acids in gels. The fluor is imbibed in the electrophoretic gel after electrophoresis; the gel is dried and then placed in contact with a preflashed film. (iii) Intensifying screens Intensifying screens consist of a flexible polyester base that is coated with inorganic phosphors. One very popular type of screen is scheelite that is calcium tungstate (CaWO4). High-energy radiation passes through the film, causes phosphor to fluoresce, and emits light, which in turn superimposes its image on the film. The intensifying screen consists of a solid phosphor, and it is placed on the other side of the film from the sample. However, the resolution of the bands is adversely affected due to the spreading of bands. This technique is very useful. It is used in gel electrophoresis or analysis of membrane filters where highenergy b-emitters (32P-labelled DNA) or g-emitting isotopes (125I-labelled protein) are used. (iv) Low-temperature exposure The exposure should be done at low temperature (about – 70°C). This is because exposure at low temperature (– 70°C) increases the sensitivity of the radiographic film to low intensities of light. (v) Pref lashing The response of a photographic emulsion to radiation is not linear and usually involves a slow initial phase (lag phase) followed by a linear phase. Sensitivity of films may be increased by preflashing. This involves an ms light flash prior to the sample being brought into juxtaposition with the film and is often used where high sensitivity is required or if results are to be quantified. (vi) Quantification In many experiments, it is important to locate than to quantify radioactivity. Though, it is possible to obtain quantitative data directly from autoradiographs by using digital image analysis, quantification is not reliable at low- or high-levels of exposure because of the lag phase (see preflashing) or saturation, respectively. Preflashing combined with fluorography or intensifying screens create the best conditions for quantitative working.
10.8 SAFETY ISSUES AND RADIO-WASTE MANAGEMENT The overall objective of radiation protection is to provide an appropriate standard of protection for man without unduly limiting the beneficial practices
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giving rise to radiation exposure. Protection from radiation is an integral component of the working infrastructure of any radiology department. The main principles of radiation protection are—to provide adequate protection from undue exposure of radiation to personnel directly or indirectly involved with radiation, without unduly limiting the benefits of radiation exposure. The components of radiation protection include justification of the procedure involving the radiation exposure, use of minimum radiation exposure compatible with the procedure which provides adequate diagnostic information, shielding of the personnel and patient from unwanted radiation exposures, and monitoring of radiation exposure to the occupational workers and the working environment. Regular surveillance of the department for radiation levels, monitoring of the radiation protection programmes, and regular educational activities form integral parts of the responsibilities of the radiation safety officer (RSO) and other administrative authorities of the department or hospital. The norms by International Commission for Radiation Protection (ICRP) and Atomic Energy Regulatory Board (AERB) have to be followed in these surveys and protection programmes. (i) Declared pregnant workers and minors Because of the increased health risks to the rapidly developing embryo and foetus, pregnant women need to avoid radiation during the entire pregnancy. The radiation dose allowed is around 0.5 rem during the entire period. This is 10 per cent of the dose limit that normally applies to radiation workers. Persons under the age of 18 years are also limited to 0.5 rem/year. (ii) Controlling radiation exposure The three basic ways of controlling exposure to harmful radiation are: (a) limiting the time spent near a source of radiation, (b) increasing the distance away from the source, and (c) using shielding to stop or reduce the level of radiation. (iii) Time The radiation dose is directly proportional to the time spent in the radiation. Therefore, a person should not stay near a source of radiation any longer than necessary. If a survey meter reads 4 mR/h at a particular location, a total dose of 4 mR will be received if a person remains at that location for one hour. In a two hour span of time, a dose of 8 mR would be received. Equation 10.17 can be used for making a simple calculation to determine the dose that will be or has been received in a radiation area. Dose = Dose rate × Time 10.17 (iv) Distance Increased distance from the source of radiation will reduce the amount of radiation received. As radiation travels from the source, it spreads out becoming less intense. This is analogous to standing near a fire. The closer a person stands to the fire, the more intense the heat feels from
322 Bioanalytical Techniques the fire. This phenomenon can be expressed by an equation known as the inverse square law, which states that as the radiation travels out from the source, the dosage decreases inversely with the square of the distance. 1 10.18 Amount of radiation μ _________2 Distance (v) Shielding Third way of reducing exposure to radiation is to place something between the radiographer and the source of radiation. In general, the more dense the material, the more shielding it will provide. The most-effective shielding is provided by depleted uranium metal. It is used primarily in g ray cameras. Depleted uranium and other heavy metals, such as tungsten, are very effective in shielding radiation because their tightly packed atoms make it hard for radiation to move through the material without interacting with the atoms. Lead and concrete are the most commonly used radiation shielding materials primarily because they are easy to work with and are readily available materials. Concrete is commonly used in the construction of radiation vaults. Some vaults will also be lined with lead sheeting to help reduce the radiation to acceptable levels on the outside. In research, the investigator should critically analyse the need of the radioisotope for any particular application. Protective clothing, gloves and glasses have to be worn. The radioactive material needs to be kept safe, secure, and well-labelled. All care to avoid spillages needs to be exercised. The working area needs to be monitored frequently for radioactivity. Examples of various tissues and their relative radio sensitivities are listed in Table 10.3. (vi) Radioactive waste management Solid waste disposal, of waste such as municipal garbage, is based on three well-known methods, namely landfills, incineration, and recycling. Underground engineered trenches Table 10.3 Radiosensitivities of different tissues Level of radiosensitivity
Example
High radiosensitivity
Lymphoid organs, bone marrow, blood, testes, ovaries, intestines
Fairly high radiosensitivity
Skin and other organs with epithelial cell lining (cornea, oral cavity, esophagus, rectum, bladder, vagina, uterine cervix, ureters)
Moderate radiosensitivity
Optic lens, stomach, growing cartilage, fine vasculature, growing bone
Fairly low radiosensitivity
Mature cartilage or bones, salivary glands, respiratory organs, kidneys, liver, pancreas, thyroid, adrenal and pituitary glands
Low radiosensitivity
Muscle, brain, spinal cord
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in near-surface disposal facilities are utilized for disposal of solid waste; these disposal sites are under continuous surveillance and monitoring. High efficiency particulate air (HEPA) filters are used for minimizing air-borne radioactivity. Sophisticated methods of landfills have also been adapted for radioactive waste. Incineration of radioactive waste needs to be handled with care such as fly ash, noxious gases, and chemical contaminants. Fine particulate filters and the gaseous effluents after incineration are first diluted and released. Recycling is done whenever feasible whereby useful radioactive elements are recovered for reuse. Solar evaporation of liquid waste, reverse osmosis, and immobilization using cement matrix are adopted depending on the form of waste.
11 Immunochemical Techniques
11.1
INTRODUCTION Modern biotechnology researchers have been constantly searching for highly specific and sensitive strategies for identification, localization, and measurement of key biochemicals. The immune system is an excellent example of such a specific and sensitive system which has evolved over a period of time for the specific recognition and selective destruction of harmful microorganisms and their toxins that invade our body. In the development of modern immunochemical techniques, the researchers have emulated the components of immune system and have increasingly employed them as probes for identification, affinity ligands, and also as therapeutic agents.
11.1.1 A Brief Introduction to Immune System The immune system has been divided into innate and adaptive immune system. The innate immune system is non-specific and forms the first line of defence against invading foreign particles or antigens. The adaptive immune system is highly specific and is activated by recognition of “non-self” antigens on the invading particles. When the immune system is exposed to the antigen for the first time, adaptive immunity against that antigen is developed. This process takes time; however, during this first exposure memory cells are formed. The memory cells are responsible for a swift immune response on re-exposure of the immune system with the same antigen. The cells of the adaptive immune system comprise mononuclear leukocytes or lymphocytes categorized into B lymphocytes (B cells) and T lymphocytes (T cells). These cells are found in lymphoid organs including bone marrow, thymus, spleen, lymph nodes, and blood. Both B cells and T cells bear cell
326 Bioanalytical Techniques surface receptors that recognize antigens specifically. T cells recognize antigens only after the antigens have been processed and presented in association with a molecule named major histocompatibility complex (MHC). Therefore, exogenous antigens are first taken up by antigen-presenting cells (APCs) which include phagocytic cells like dendritic cells, macrophages, and B cells. The APCs engulf the antigen by endocytosis. The antigen is degraded into peptides. These antigenic peptides are then displayed at the surface of the cell in association with histocompatibility molecule. There are two major types of T cells: T cytotoxic (TC) cell also known as cytotoxic killer cells and T helper (T H ) cell. The T helper cells can be identified as the cell population that bears CD4 cell surface marker whereas cytotoxic killer cells possess CD8 cell surface marker. T cell receptor (TCR) of TC cells recognize antigens coupled to Class I MHC molecules, while T H cells recognize antigens coupled to Class II MHC molecules. T H cells release cytokines when they encounter cells that have antigen bound to MHC Class II molecules resulting in killing of host cells. Cytokines are proteins produced by T H cells that enhance the activity of macrophages, TH cells, and antibodyproducing B cells. T cell activation is closely controlled and requires a very strong MHC–antigen activation signal along with additional activation signals provided by cytokines. T H cells have no cytotoxic activity but they control the immune response by affecting the response from other cells. The B-cell receptor (BCR) is an antibody molecule bound on the surface of B cells. It recognizes the antigens without the requirement of antigen processing. Each lineage of B cell expresses a different antibody that recognizes and binds to a specific antigen. This antigen–antibody complex is taken up by the B cell and processed by proteolytic enzymes into peptides. The B cell then displays these antigenic peptides on its surface bound to MHC Class II molecules. The MHC–antigen complex attracts a corresponding T H cell, which on interaction releases lymphokines and activates the B cells. The activated B cells then begin to differentiate into plasma cells which secrete millions of copies of the antibody against this antigen (Figure 11.1). These antibodies circulate in blood plasma and lymph, bind to pathogens expressing the antigen, and mark them for destruction by complement activation or for uptake and destruction by phagocytes. Antibodies can also neutralize the foreign invasion directly by binding to bacterial toxins or by interfering with the receptors that virus and bacteria use for infecting cells. During first exposure of antigen, the memory cells are formed that mount a prompt immune response on subsequent exposure to the same antigen.
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Figure 11.1
327
Components of an immune system
11.1.2 Antigens Antigens are the foreign molecules that are recognized and bound by the antibodies. They are most commonly polypeptides or carbohydrates in nature, but they can also be lipids, nucleic acids, or even small molecules such as neurotransmitters. An antibody molecule interacts with only a small region or portion of an antigen. The region of an antigen recognized by an antibody is called an epitope. Most antigens are complex proteins containing different antigenic determinants, and the immune system usually responds by producing antibodies to several epitopes on the antigen. The antigenic determinants can be continuous or discontinuous depending upon their distribution in different regions of a primary structure of the protein. Discontinuous determinants are brought together due to the formation of the secondary or tertiary structure of the protein. Complete antigens are molecules that can induce specific immune responses and react with the products of those responses. An incomplete antigen is called a hapten and is a chemically defined substance of low molecular weight that cannot induce an immune response on its own but can react with the specific antibodies. A wide variety of molecules, including drugs, explosives, pesticides, herbicides, polycyclic aromatic hydrocarbons, and metal ions can be considered as haptens. These haptens can induce the immune system to produce antibodies only when they are covalently conjugated to a larger carrier molecule such as a protein.
328 Bioanalytical Techniques 11.1.3 Antibody Electrophoresis of serum revealed bands corresponding to albumin and the alpha (a), beta (b), and gamma (g) globulins. When the g-globulin fraction of the serum was separated into high- and low-molecular weight fractions, antibodies having molecular weight of about 150 kD, nominated as immunoglobulin G (IgG), were found in the low-molecular weight fraction. This fraction was found to interact with specific antigens. Therefore, the g-globulin fraction was established to contain serum antibodies. The antibodies were found to be glycoproteins which were designated as immunoglobulins, to distinguish them from any other proteins that might be contained in the g-globulin fraction. Antibody molecules are polypeptides consisting of four peptide chains (Figure 11.2). The structure consists of two identical light (L) chains, polypeptides of about 25 kD (about 214 amino acids) and two identical heavy (H) chains, larger polypeptides of molecular weight 50 kD or more (about 446 amino acids). Each L chain is bound to a H chain by a disulfide bond, and by non-covalent interactions such as salt linkages, hydrogen bonds, and hydrophobic bonds, to form a heterodimer (H–L). Similar non-covalent interactions and disulfide
Figure 11.2
Structure of an immunoglobulin
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bridges link the two identical H–L chain combinations to each other to form the basic four-chain structure. The exact number and positions of disulfide bonds differ among antibody classes and subclasses. The first 110 amino acids or so of the amino-terminal region of a L or H chain vary greatly among antibodies of different specificity. These segments of highly variable sequence are called V regions: V L in light chains and V H in heavy chains. The difference in specificity displayed by different antibodies is the result of differences in the amino acid sequences of variable regions. These regions are called as complementaritydetermining regions (CDRs), and it is these CDRs, on both light and heavy chains, that constitute the antigen-binding site of the antibody molecule. The regions of relatively constant sequence beyond the variable regions have been named as C regions, CL on the light chain and CH on the heavy chain. The constant region of the antibody also serves as a useful handle to manipulate the antibody during most of the immunochemical techniques. The five primary classes of immunoglobulins are IgG, IgM, IgA, IgD, and IgE. These are distinguished by the type of H chain found in the molecule such as gamma (g) chains in IgG, mu (m) chains in IgM, alpha (a) chains in IgA, epsilon (e) chains in IgE, and delta (d) chains in IgD. There are only two main types of L chains: kappa (k) and lambda (l). The number of antigen-binding variable regions on the antibody corresponds with its subclass, and determines the valence of the antibody. For example, in humans, functioning IgM antibodies have five Y-shaped units (pentamer) containing a total of 10 light chains, 10 heavy chains, and 10 antigen-binding sites. The most commonly used antibody in immunochemical procedures is of the IgG class since it is the most abundant of the antibodies.
11.2
PRODUCTION AND PURIFICATION OF ANTIBODIES
11.2.1
Production of Polyclonal Antibodies (Antisera)
Polyclonal antibodies are secreted by different lineages of B cell within the body. They are a collection of immunoglobulin molecules that react against a specific antigen, each identifying a different epitope. Polyclonal antibody production is carried out using laboratory animals such as rabbits. When selecting the animal species for polyclonal antibody production, it is important to consider the following: (1) the amount of polyclonal antibody needed, (2) the ease of obtaining blood samples, (3) the phylogenetic relationship between the antigen and the animal species, and (4) the intended use of the polyclonal antibody. The most frequently used animal species for polyclonal antibody production are rabbit, mouse, rat, hamster, guinea pig, goat, sheep, and chicken.
330 Bioanalytical Techniques Polyclonal antibodies are produced in animals by immunizing them with antigen (also known as immunogen) prepared with an adjuvant. Adjuvants are chemicals such as detergents and oils or complex proprietary products containing bacterial cell walls or their preparations which can boost the immune response. For smaller or less immunogenic proteins or peptides, the immunogen can also be coupled to carrier proteins such as keyhole limpet hemocyanin (KLH), bovine serum albumin (BSA), ovalbumin (OVA), and purified protein derivative (PPD) of tuberculin. Immunization is typically performed by injecting the antigen formulated with the adjuvant intra-dermally or subcutaneously. Repeated immunizations in the model organism is carried out at intervals of about 4–6 weeks. During this immunization schedule, the immunized animals are regularly monitored for the titres of antibody in their blood samples (Figure 11.3). Once high level of circulating antibody is detected in test bleeds, then the blood harvesting is carried out to obtain the antibody. Blood is harvested until the antibody titre begins to fall. A second round of blood can also be harvested from the animal after repeat immunization. The harvested blood is allowed to clot and the serum is collected and stored at 4°C or lower. Polyclonal antibodies will recognize multiple epitopes on an antigen. Therefore, polyclonal antibodies identify proteins of high homology to the immunogen protein and these can be used for screening the target protein in
Figure 11.3
Process of polyclonal antibody production
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samples where nature of the antigen is not known. Polyclonal antibodies can amplify signal from target protein as they can bind to different epitopes in an antigen. Polyclonal antibodies can react with respective antigen even if the antigen itself is altered due to polymorphism, heterogeneity of glycosylation, or slight denaturation, as compared to monoclonal (homogeneous) antibodies. Polyclonal antibodies are often the preferred choice for detection of denatured proteins since recognition of multiple epitopes aids the detection. Due to recognition of multiple epitopes, polyclonal antibodies can give better results in immunoprecipitation. Despite all the advantages, polyclonal antibodies suffer from batch to batch variability and non-specificity leading to background noise. Polyclonal antibodies are not useful for probing specific domains of antigen, because antiserum will usually recognize many domains. The incidence of cross-reactivity due to recognition of multiple epitopes needs to be assessed. Moreover, it would not be suitable for quantification experiments, for, example in flow cytometry, as the results would become inaccurate.
11.2.2
Production of Monoclonal Antibodies
A monoclonal antibody is produced by a single clone of B lymphoid cells. All of the molecules of a given monoclonal antibody have the same amino acid sequence and, hence, the same binding properties. Use of monoclonal antibodies rather than the traditional polyclonal antibodies eliminates the difficulties like batch to batch variation, non-specificity, and cross-reactivity. The antigen–monoclonal antibody reactions are quantitative and therefore results of the tests posses the desired accuracy for use in diagnosis and research. (i) Hybridoma process for monoclonal antibody production In 1975, J. F. Köhler and César Milstein made a pioneering development of a technology that made the production of mono-specific antibodies called monoclonal antibodies. They used a hybridoma cell, a hybrid cell formed by the fusion of plasma cells producing specific antibody and a myeloma cell (a malignant or cancerous cell). The hybridoma cell possesses the qualities of both the parent cells, that is the ability to grow continually and to produce large amounts of pure antibodies. They have high growth rate of the myeloma cell and, hence, large-scale culture and production are uncomplicated. For the production of monoclonal antibodies, laboratory mice are first exposed to an antigen for the production of the desired antibody. The specific B lymphocytes are isolated from the spleen of the immunized mice (Figure 11.4).
332 Bioanalytical Techniques
Rat immunized with the desired antigen
B lymphocyted isolated from rat spleen
Myeloma cell culture
B lymphocyted fused with myeloma cells in fusion medium containing PEG or Sendai virus
The cells seeded in HAT selection medium
Only hybridoma cells survive in the HAT medium
Hybridoma cells screened for the production of desired antibody Clonal expansion of hybridoma cells
Figure 11.4
Hybridoma technique for the production of monoclonal antibody
The isolated B lymphocytes are fused with immortalized myeloma cells in a fusion medium using polyethylene glycol (PEG) or Sendai virus. Subsequently, the cells are incubated in a selective medium known as hypoxanthine aminopterin thymidine (HAT) medium. HAT selection depends on the fact that mammalian cells can synthesize nucleotides by two different pathways: the de novo and salvage pathways (Figure 11.5). The de novo pathway is blocked by aminopterin, a folic acid analogue. When the de novo pathway is blocked, cells utilize the salvage pathway for nucleotide synthesis. The enzymes catalysing the salvage pathway include hypoxanthine-guanine phosphorribosyl transferase (HGPRT) and thymidine kinase (TK). A mutation in either of these two enzymes blocks the salvage pathway. HAT medium contains aminopterin to block the de novo pathway, and hypoxanthine and thymidine to allow growth via the
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Figure 11.5
333
Nucleotide synthesis pathways in mammalian cells
Note: Aminopterin in the medium blocks the De novo pathway for nucleotide synthesis. Myeloma cells carry a mutation in either Thymidine Kinase (TK) or Hypoxanthine Guanine phosphoribosyl transferase (HGPRT) genes
salvage pathway. The myeloma cells used in this procedure carry a mutation in either TK or HGPRT gene, and, hence, these enzymes are not functional in these cells. Therefore, in the HAT medium only the hybrid cells can survive as they possess the necessary enzymes for growth on HAT medium. B lymphocytes have short lifespan; therefore, even though they possess the machinery for de novo pathway, they eventually die in culture. The antibody is secreted by the hybridoma cells in the culture medium. The culture medium is first conditioned or prepared for purification. Cells, cell debris, lipids, and clotted material are first removed, typically by centrifugation followed by filtration with a 0.45 µm filter. These large particles can cause a phenomenon called membrane fouling in later purification steps. In addition, the concentration of product in the sample may not be sufficient, especially in the cases where the desired antibody is produced by a low-secreting cell line. The sample is therefore condensed by ultrafiltration or dialysis. Most of the charged impurities are usually anions such as nucleic acids and endotoxins. These are often separated by ion-exchange chromatography. Either cation-exchange chromatography is used at such a low pH that the desired antibody binds to the column while anions flow through, or anion-exchange chromatography is used at such a high pH that the desired antibody flows through the column while anions bind
334 Bioanalytical Techniques to it. Various proteins can also be separated out along with the anions based on their isoelectric point (pI). For example, albumin has a pI of 4.8 that is significantly lower than most of the monoclonal antibodies, which have a pI of 6.1. In other words, at a given pH, the average charge of albumin molecules is likely to be more negative. Transferrin, on the other hand, has a pI of 5.9, so it cannot easily be separated using ion-exchange chromatography as pI should differ at least by 1 for a good separation. Transferrin can instead be removed by size exclusion chromatography. The advantage of this purification method is that it is one of the more reliable chromatographic techniques. Since we are dealing with proteins, properties such as charge and affinity vary with pH as molecules are protonated and deprotonated, while the size stays relatively constant. Nonetheless, it has drawbacks such as low resolution, low capacity, and low elution times. A much quicker, single-step method of separation is protein A/G affinity chromatography. The antibody selectively binds to protein A or G, so a high level of purity (generally more than 80 per cent) is obtained. However, this method may be unsuitable for antibodies that are easily damaged, as harsh conditions are generally used. A low pH can break the bonds between the antibody and the protein eliminating it from the column. In addition to possibly affecting the product, low pH can cause protein A and G itself to leak off the column and appear in the eluted sample. Gentle elution buffer systems that employ high salt concentrations are also available to avoid exposing sensitive antibodies to low pH. Cost is also an important consideration with this method because immobilized protein A or G is a more expensive resin. (ii) Antibody production in ascites f luid Small quantities of very pure polyclonal and monoclonal antibodies can be produced in rats and mice in ascites fluid. When a tumour is established in the peritoneum (cavity containing the intestines), fluid similar to plasma is secreted into the cavity of the animal in large quantities. Animals are first challenged with the antigen of interest and once a high serum level of the antibody is detected then ascites fluid production is induced. Non-secretory myeloma cells such as NSO cells are introduced into the peritoneal cavity of the animal by injection and allowed to grow there. The presence of the tumour cells causes the animal to produce ascites fluid which contains high levels of immunoglobulins to the original antigen. The fluid can be harvested by aspiration with a syringe and the antibodies can be purified. In this way polyclonal antibodies can be produced in mice. Alternatively, hybridoma cells can be injected into the peritoneal cavity and allowed to grow there.
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The peritoneum serves the purpose of an incubator that allows the growth of the hybridoma cells and monoclonal antibody can be harvested easily.
11.2.3
Chicken Egg Antibodies
The production of chicken egg antibody known as immunoglobulin Y (IgY) in chicken eggs is a simple, non-invasive, and cost-effective technique for the production of antibodies. Chicken egg antibody is the major antibody in birds and reptiles. Chickens secrete it into eggs to provide protection for the developing embryo. The production process starts by immunization of the chickens with the antigen. The process of immunization is repeated three or four times and the immune status is monitored by test bleeds. The eggs that are collected may contain up to 50 mg of antibody per yolk. A number of methods have been standardized for the purification of the immunoglobulin Y antibody from the egg yolks. It has been observed that antigens that give a poor immunogenic response in mammals can give much higher yields in chickens. Besides these there are several other advantages of using IgY. Repeated egg laying leads to high antibody yields per hen. The IgY has markedly less chances of cross-reactivity with mammalian antibodies as compared to that of IgG. Of the immunoglobulins arising during the immune response, only IgY is found in chicken eggs. Thus, in preparations from chicken eggs, there is no contamination with IgA or IgM. The yield of IgY from a chicken egg is comparable to that of IgG from rabbit serum. However, the isolation of IgY from egg yolk is more difficult than the isolation of IgG from blood serum as it does not bind to protein A or G. Additionally, the egg yolk contains high concentration of lipids which need to be removed during purification.
11.2.4 Applications of Monoclonal Antibodies Monoclonal antibodies have a variety of academic, medical, and commercial uses. The prominent applications have been listed below. However, the list is not exhaustive but indicative of the importance of monoclonal antibodies. (i) Antibodies are used in several diagnostic tests to detect small amounts of drugs, toxins, or hormones. For example, antibodies are used for the diagnosis of AIDS by the enzyme-linked immunosorbent assay (ELISA) test and monoclonal antibodies to human chorionic gonadotropin (HCG) are used in pregnancy test kits. (ii) Antibodies are used in the radioimmunoassay and radioimmunotherapy of cancer. (iii) Monoclonal antibodies can be used for treating viral diseases such as AIDS.
336 Bioanalytical Techniques (iv) Monoclonal antibodies can be used for classifying strains of a single pathogen; for example, Neisseria gonorrhoeae can be typed using monoclonal antibodies. (v) Researchers use monoclonal antibodies to identify and trace specific cells or molecules in an organism. (vi) An antibody to the T3 antigen of T cells has been used for lessening the problem of organ rejection in patients who have had organ transplants. The advantages of using monoclonal antibodies as compared to polyclonal antibodies are as follows: • Higher specificity makes monoclonal antibodies less likely to cross react with other proteins. • These are excellent as the primary antibody in an assay, or for detecting antigens in tissue, since they give significantly less background staining. • If experimental conditions are kept constant, results from monoclonal antibodies will be highly reproducible. • They are extremely efficient for binding of antigen within a mixture of related molecules, such as in the case of affinity purification. The disadvantages are as follows: • They are less likely to detect the antigen across a range of species. • It may not interact with the antigen even if there is a very slight change in the structure. This can be overcome by pooling two or more monoclonal antibodies to the same antigen for such an application.
11.3
IMMUNOASSAY TECHNIQUES An immunoassay is a biochemical test that measures the presence or concentration of a macromolecule in a solution through the use of an antibody or immunoglobulin. Immunoassay is the method of choice for measuring low concentration of analytes that cannot be detected by other methods. Common uses include measurement of drugs, hormones, proteins, tumour markers, and markers of cardiac injury. Very sensitive and specific immunoassays are available for a number of diagnostic as well as research applications (Table 11.1). Qualitative immunoassays are often used for detecting antigens on infectious agents and antibodies that the body produces to fight them. For example, immunoassays are used for detecting antigens on Hemophilus, Cryptococcus, and Streptococcus organisms in the cerebrospinal fluid (CSF) of meningitis patients. They are also used for detecting antigens associated with organisms that are difficult to culture, such as hepatitis B virus (HBV) and
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Table 11.1 Types of immunoassays and their sensitivities Assay Precipitation reaction in fluids Radial immunodiffusion
Sensitivity (mg Antibody/ml) 20–200 10–50
Ouchterlony double immunodiffusion
20–200
Agglutination: Direct
20–200
Agglutination: Passive
0.3
Agglutination inhibition
0.0006–0.006
ELISA